 • Wikijunior •  •  • Using Wikibooks  to , the open-content textbooks collection that . books with . 

Welcome to the world's foremost open content&lt;br&gt;Organic Chemistry Textbook&lt;br&gt;on the web! The Study of Organic Chemistry. Organic chemistry is primarily devoted to the unique properties of the carbon atom and its compounds. These compounds play a critical role in biology and ecology, Earth sciences and geology, physics, industry, medicine and — of course — chemistry. At first glance, the new material that organic chemistry brings to the table may seem complicated and daunting, but all it takes is concentration and perseverance. Millions of students before you have successfully passed this course and you can too! This field of chemistry is based less on formulas and more on reactions between various molecules under different conditions. Whereas a typical general chemistry question may ask a student to compute an answer with an equation from the chapter that they memorized, a more typical organic chemistry question is along the lines of "what product will form when substance X is treated with solution Y and bombarded by light". The key to learning organic chemistry is to "understand" it rather than cram it in the night before a test. It is all well and good to memorize the mechanism of Michael addition, but a superior accomplishment would be the ability to explain "why" such a reaction would take place. As in all things, it is easier to build up a body of new knowledge on a foundation of solid prior knowledge. Students will be well served by much of the knowledge brought to this subject from the subject of General Chemistry. Concepts with particular importance to organic chemists are covalent bonding, Molecular Orbit theory, VSEPR Modeling, understanding acid/base chemistry vis-a-vis pKa values, and even trends of the periodic table. This is by no means a comprehensive list of the knowledge you should have gained already in order to fully understand the subject of organic chemistry, but it should give you some idea of the things you need to know to succeed in an organic chemistry test or course. Organic Chemistry is one of the subjects which are very useful and close to our daily life. We always try to figure out some of the unknown mysteries of our daily life through our factious thinking habit, which generates superstitions. Through the help of chemistry we can help ourselves to get out of this kind of superstition. We always try to find the ultimate truth through our own convenience. In the ancient past we had struggled to make things to go as per our need. In that context we have found fire, house, food, transportation, etc... Now the burning question is: "how can chemistry help our daily life?" To find the answer of this questions, we have to know the subject thoroughly. Let us start it from now. 

The purpose of this section is to review topics from freshman chemistry and build the foundation for further studies in organic chemistry.  | Alkanes » 

 &gt; Introduction to reactions 

« Haloalkanes |Alkenes| Alkynes » Alkenes are aliphatic hydrocarbons containing carbon-carbon double bonds and general formula CnH2n. =Naming Alkenes= Alkenes are named as if they were alkanes, but the "-ane" suffix is changed to "-ene". If the alkene contains only one double bond and that double bond is terminal (the double bond is at one end of the molecule or another) then it is not necessary to place any number in front of the name. butane: C4H10 (CH3CH2CH2CH3)&lt;br&gt; butene: C4H8 (CH2=CHCH2CH3) If the double bond is not terminal (if it is on a carbon somewhere in the center of the chain) then the carbons should be numbered in such a way as to give the first of the two double-bonded carbons the lowest possible number, and that number should precede the "ene" suffix with a dash, as shown below. correct: pent-2-ene (CH3CH=CHCH2CH3)&lt;br&gt; incorrect: pent-3-ene (CH3CH2CH=CHCH3)&lt;br&gt; "The second one is incorrect because flipping the formula horizontally results in a lower number for the alkene." If there is more than one double bond in an alkene, all of the bonds should be numbered in the name of the molecule - even terminal double bonds. The numbers should go from lowest to highest, and be separated from one another by a comma. The IUPAC numerical prefixes are used to indicate the number of double bonds. octa-2,4-diene: CH3CH=CHCH=CHCH2CH2CH3&lt;br&gt; deca-1,5-diene: CH2=CHCH2CH2CH=CHCH2CH2CH2CH3 Note that the numbering of "2-4" above yields a molecule with two double bonds separated by just one single bond. Double bonds in such a condition are called "conjugated", and they represent an enhanced stability of conformation, so they are energetically favored as reactants in many situations and combinations. EZ Notation. Earlier in stereochemistry, we discussed cis/trans notation where cis- means same side and trans- means opposite side. Alkenes can present a unique problem, however in that the cis/trans notation sometimes breaks down. The first thing to keep in mind is that alkenes are planar and there's no rotation of the bonds, as we'll discuss later. So when a substituent is on one side of the double-bond, it stays on that side. The above example is pretty straight-forward. On the left, we have two methyl groups on the same side, so it's cis-but-2-ene. And on the right, we have them on opposite sides, so we have trans-but-2-ene. So in this situation, the cis/trans notation works and, in fact, these are the correct names. From the example above, how would you use cis and trans? Which is the same side and which is the opposite side? Whenever an alkene has 3 or 4 differing substituents, one must use the what's called the EZ nomenclature, coming from the German words, Entgegen (opposite) and Zusammen (same). Let's begin with (Z)-3-methylpent-2-ene. We begin by dividing our alkene into left and right halves. On each side, we assign a substituent as being either a high priority or low priority substituent. The priority is based on the atomic number of the substituents. So on the left side, hydrogen is the lowest priority because its atomic number is 1 and carbon is higher because its atomic number is 6. On the right side, we have carbon substituents on both the top and bottom, so we go out to the next bond. On to the top, there's another carbon, but on the bottom, a hydrogen. So the top gets high priority and the bottom gets low priority. Because the high priorities from both sides are on the same side, they are Zusammen (as a mnemonic, think 'Zame Zide'). Now let's look at (E)-3-methylpent-2-ene. On the left, we have the same substituents on the same sides, so the priorities are the same as in the Zusammen version. However, the substituents are reversed on the right side with the high priority substituent on the bottom and the low priority substituent on the top. Because the High and Low priorities are opposite on the left and right, these are Entgegen, or opposite. The system takes a little getting used to and it's usually easier to name an alkene than it is to write one out given its name. But with a little practice, you'll find that it's quite easy. Comparison of E-Z with cis-trans. To a certain extent, the Z configuration can be regarded as the "cis-" isomer and the E as the "trans-" isomers. This correspondence is exact only if the two carbon atoms are identically substituted. In general, cis-trans should only be used if each double-bonded carbon atom has a hydrogen atom (i.e. R-CH=CH-R'). IUPAC Gold book on cis-trans notation. IUPAC Gold book on E-Z notation. =Properties= Alkenes are molecules with carbons bonded to hydrogens which contain at least two sp2 hybridized carbon atoms. That is, to say, at least one carbon-to-carbon double bond, where the carbon atoms, in addition to an electron pair shared in a "sigma" (σ) bond, share one pair of electrons in a "pi" (π) bond between them. The general formula for an aliphatic alkene is: CnH2n -- "e.g." C2H4 "or" C3H6 Diastereomerism. Restricted rotation. Because of the characteristics of pi-bonds, alkenes have very limited rotation around the double bonds between two atoms. In order for the alkene structure to rotate the pi-bond would first have to be broken - which would require about 60 or 70 kcal of energy per mol. For this reason alkenes have different chemical properties based on which side of the bond each atom is located. For example, but-2-ene exists as two diastereomers: =Relative stability= Observing the reaction of the addition of hydrogen to 1-butene, (Z)-2-butene, and (E)-2-butene, we can see that all of the products are butane. The difference between the reactions is that each reaction has a different energy: -30.3 kcal/mol for 1-butene, -28.6 kcal/mol for (Z)-2-butene and -27.6 kcal/mol for (E)-2-butene. This illustrates that there are differences in the stabilities of the three species of butene isomers, due to the difference in how much energy can be released by reducing them. The relative stability of alkenes may be estimated based on the following concepts: Internal alkenes are more stable than terminal alkenes because they are connected to more carbons on the chain. Since a terminal alkene is located at the end of the chain, the double bond is only connected to one carbon, and is called primary (1°). Primary carbons are the least stable. In the middle of a chain, a double bond could be connected to two carbons. This is called secondary (2°). The most stable would be quaternary (4°). =Reactions= Preparation. There are several methods for creating alkenes. Some of these methods, such as the Wittig reaction, we'll only describe briefly in this chapter and instead, cover them in more detail later in the book. For now, it's enough to know that they are ways of creating alkenes. Dehydrohalogenation of Haloalkanes. Alkyl halides are converted into alkenes by dehydrohalogenation: elimination of the elements of hydrogen halide. Dehydrohalogenation involves removal of the halogen atom together with a hydrogen atom from a carbon adjacent to the one bearing the halogen. It uses the E2 elimination mechanism that we'll discuss in detail at the end of this chapter The haloalkane must have a hydrogen and halide 180° from each other on neighboring carbons. If there is no hydrogen 180° from the halogen on a neighboring carbon, the reaction will not take place. It is not surprising that the reagent required for the elimination of what amounts to a molecule of acid is a strong base for example: alcholic KOH. In some cases this reaction yields a single alkene. and in other cases yield a mixture. n-Butyl chloride, for example, can eliminate hydrogen only from C-2 and hence yields only 1-butene. sec-Butyl chloride, on the other hand, can eliminate hydrogen from either C-l or C-3 and hence yields both 1-butene and 2-butene. Where the two alkenes can be formed, 2-butene is the chief product. Dehalogenation of Vicinal Dihalides. The dehalogenation of vicinal dihalides (halides on two neighboring carbons, think "vicinity") is another method for synthesizing alkenes. The reaction can take place using either sodium iodide in a solution of acetone, or it can be performed using zinc dust in a solution of either heated ethanol or acetic acid. This reaction can also be performed with magnesium in ether, though the mechanism is different as this actually produces, as an intermediate, a Grignard reagent that reacts with itself and causes an elimination, resulting in the alkene. Dehydration of alcohols. An alcohol is converted into an alkene by dehydration: elimination of a molecule of water. Dehydration requires the presence of an acid and the application of heat. It is generally carried out in either of two ways, heating the alcohol with sulfuric or phosphoric acid to temperatures as high as 200, or passing the alcohol vapor over alumina, Al2O3 , at 350-400, alumina here serving as a Lewis acid. Ease of dehydration of alcohols : 3° &gt; 2° &gt; 1° Where isomeric alkenes can be formed, we again find the tendency for one isomer to predominate. Thus, sec-butyl alcohol, which might yield both 2-butene and 1-butene, actually yields almost exclusively the 2-isomer The formation of 2-butene from n-butyl alcohol illustrates a characteristic of dehydration that is not shared by dehydrohalogenalion: the double bond can be formed at a position remote from the carbon originally holding the -OH group. This characteristic is accounted for later. It is chiefly because of the greater certainty as to where the double bond will appear that dehydrohalogeation is often preferred over dehydration as a method of making alkenes. Reduction of Alkynes. Reduction of an alkyne to the double-bond stage can yield either a cis-alkene or a trans-alkene, unless the triple bond is at the end of a chain. Just which isomer predominates depends upon the choice of reducing agent. Predominantly trans-alkene is obtained by reduction of alkynes with sodium or lithium in liquid ammonia. Almost entirely cis-alkene (as high as 98%) is obtained by hydrogenation of alkynes with several different catalysts : a specially prepared palladium called Lindlar's catalyst; or a nickel boride called P-2 catalyst. Each of these reactions is, then, highly stereoselective. The stereoselectivity in the cis-reduction of alkynes is attributed, in a general way, to the attachment of two hydrogens to the same side of an alkyne sitting on the catalyst surface; presumably this same stereochemistry holds for the hydrogenation of terminal alkynes which cannot yield cis- and trans-alkenes. Markovnikov's Rule. Before we continue discussing reactions, we need to take a detour and discuss a subject that's very important in Alkene reactions, "Markovnikov's Rule." This is a simple rule stated by the Russian Vladmir Markovnikov in 1869, as he was showing the orientation of addition of HBr to alkenes. His rule states:"When an unsymmetrical alkene reacts with a hydrogen halide to give an alkyl halide, the hydrogen adds to the carbon of the alkene that has the greater number of hydrogen substituents, and the halogen to the carbon of the alkene with the fewer number of hydrogen substituents" (This rule is often compared to the phrase: "The rich get richer and the poor get poorer." Aka, the Carbon with the most Hydrogens gets another Hydrogen and the one with the least Hydrogens gets the halogen) This means that the nucleophile of the electophile-nucleophile pair is bonded to the position most stable for a carbocation, or partial positive charge in the case of a transition state. Examples. formula_1 Here the Br attaches to the middle carbon over the terminal carbon, because of Markovnikov's rule, and this is called a Markovnikov product. Markovnikov product. The product of a reaction that follows Markovnikov's rule is called a Markovnikov product. Markovnikov addition. Markovnikov addition is an addition reaction which follows Markovnikov's rule, producing a Markovnikov product. Anti-Markovnikov addition. Certain reactions produce the opposite of the Markovnikov product, yielding what is called anti-Markovnikov product. That is, hydrogen ends up on the more substituted carbon of the double bond. The hydroboration/oxidation reaction that we'll discuss shortly, is an example of this, as are reactions that are conducted in peroxides. A modernized version of Markovnikov's rule often explains the "anti-Markovnikov" behavior. The original Markovnikov rule predicts that the hydrogen (an electrophile) being added across a double bond will end up on the carbon with more hydrogens. Generalizing to all electrophiles, it is really the electrophile which ends up on the carbon with the greatest number of hydrogens. Usually hydrogen plays the role of the electrophile; however, hydrogen can also act as an nucleophile in some reactions. The following expansion of Markovnikov's rule is more versatile: "When an alkene undergoes electrophilic addition, the electrophile adds to the carbon with the greatest number of hydrogen substituents. The nucleophile adds to the more highly substituated carbon." Or more simply: "The species that adds first adds to the carbon with the greatest number of hydrogens." The fact that some reactions reliably produce anti-Markovnikov products is actually a powerful tool in organic chemistry. For example, in the reactions we discuss below, we'll show two different ways of creating alcohols from alkenes: Oxymercuration-Reduction and Hydroboration/Oxidation. Oxymercuration produces a Markovnikov product while Hydroboration produces an anti-Markovnikov product. This gives the organic chemist a choice in products without having to be stuck with a single product that might not be the most desired. Why it works. Markovnikov's rule works because of the stability of carbocation intermediates. Experiments tend to reveal that carbocations are planar molecules, with a carbon that has three substituents at 120° to each other and a vacant p orbital that is perpendicular to it in the 3rd plane. The p orbital extends above and below the trisubstituent plane. This leads to a stabilizing effect called hyperconjugation. Hyperconjugation is what happens when there is an unfilled (antibonding or vacant) C-C π orbital and a filled C-H σ bond orbital next to each other. The result is that the filled C-H σ orbital interacts with the unfilled C-C π orbital and stabilizes the molecule. The more highly substituted the molecule, the more chances there are for hyperconjugation and thus the more stable the molecule is. Another stabilizing effect is an inductive effect. Exceptions to the Rule. There are a few exceptions to the Markovnikov rule, and these are of tremendous importance to organic synthesis. Addition reactions. Hydroboration. Hydroboration is a very useful reaction in Alkenes, not as an end product so much as an intermediate product for further reactions. The primary one we'll discuss below is the Hydroboration/Oxidation reaction which is actually a hydroboration reaction followed by a completely separate oxidation reaction. The addition of BH3 is a concerted reaction in that several bonds are broken and formed at the same time. Hydroboration happens in what's called syn-addition because the boron and one of its hydrogens attach to the same side of the alkene at the same time. As you can see from the transition state in the center of the image, this produces a sort of box between the two alkene carbons and the boron and its hydrogen. In the final step, the boron, along with its other two hydrogens, remains attached to one carbon and the other hydrogen attaches to the adjacent carbon. This description is fairly adequate, however, the reaction actually continues to happen and the -BH2 continue to react with other alkenes giving an R2BH and then again, until you end up with a complex of the boron atom attached to 3 alkyl groups, or R3B. This trialkyl-boron complex is then used in other reactions to produce various products. Borane, in reality, is not stable as BH3. Boron, in this configuration has only 6 electrons and wants 8, so in its natural state it actually creates the B2H6 complex shown on the left. Furthermore, instead of using B2H6 itself, BH3 is often used in a complex with tetrahydrofuran (THF) as shown in the image on the right.In either situation, the result of the reactions are the same. Hydroboration/Oxidation. With the reagent diborane, (BH3)2, alkenes undergo hydroboration to yield alkylboranes, R3B, which on oxidation give alcohols.The reaction procedure is simple and convenient, the yields are exceedingly high, and the products are ones difficult to obtain from alkenes in anyother way. Diborane is the dimer of the hypothetical BH3 (borane) and, in the reactions that concern us, acts much as though it were BH3 . Indeed, in tetrahydrofuran, one of the solvents used for hydroboration, the reagent exists as the monomer, in the form of an acid-base complex with the solvent. Hydroboration involves addition to the double bond of BH3 (or, in following stages, BH2R and BHR2), with hydrogen becoming attached to one doubly-bonded carbon, and boron to the other. The alkylborane can then undergo oxidation, in which the boron is replaced by -OH. Thus, the two-stage reaction process of hydroboration-oxidation permits, in effect, the addition to the carbon-carbon double bond of the elements of H-OH. Reaction is carried out in an ether, commonly tetrahydrofuran or "diglyme" (diethylene glycol methyl ether, CH3OCH2CH2OCH2CH2OCH3). Diborane is commercially available in tetrahydrofuran solution. The alkylboranes are not isolated, but are simply treated in situ with alkaline hydrogen peroxide. Stereochemistry and Orientation. Hydroboration-oxidation, then, converts alkenes into alcohols. Addition is highly regiospecific; the preferred product here, however, is exactly opposite to the one formed by oxymercuration-demercuration or by direct acid-catalyzed hydration. The hydroboration-oxidation process gives products corresponding to anti-Markovnikov addition of water to the carbon-carbon double bond. The reaction of 3,3-dimethyl-l -butene illustrates a particular advantage of the method. Rearrangement does not occur in hydroboration evidently because carbonium ions are not intermediates and hence the method can be used without the complications that often accompany other addition reactions. The reaction of 1,2-dimethylcyclopentene illustrates the stereochemistry of the synthesis: hydroboration-oxidation involves overall syn addition. Oxymercuration/Reduction. Alkenes react with mercuric acetate in the presence of water to give hydroxymercurial compounds which on reduction yield alcohols. The first stage, oxymercuration, involves addition to the carbon-carbon double bond of -OH and -HgOAc. Then, in reduction, the -HgOAc is replaced by -H. The reaction sequence amounts to hydration of the alkene, but is much more widely applicable than direct hydration. The two-stage process of oxymercuration/reduction is fast and convenient, takes place under mild conditions, and gives excellent yields often over 90%. The alkene is added at room temperature to an aqueous solution of mercuric acetate diluted with the solvent tetrahydrofuran. Reaction is generally complete within minutes. The organomercurial compound is not isolated but is simply reduced in situ by sodium borohydride, NaBH4. (The mercury is recovered as a ball of elemental mercury.) Oxymercuration/reduction is highly regiospecific, and gives alcohols corresponding to Markovnikov addition of water to the carbon-carbon double bond. Oxymercuration involves electrophilic addition to the carbon-carbon double bond, with the mercuric ion acting as electrophile. The absence of rearrangement and the high degree of stereospecificity (typically anti) in the oxymercuration step argues against an open carbonium ion as intermediate. Instead, it has been proposed, there is formed a cyclic mercurinium ion, analogous to the bromonium and chloronium ions involved in the addition of halogens. In 1971, Olah reported spectroscopic evidence for the preparation of stable solutions of such mercurinium ions. The mercurinium ion is attacked by the nucleophilic solvent water, in the present case to yield the addition product. This attack is back-side (unless prevented by some structural feature) and the net result is anti addition, as in the addition of halogens. Attack is thus of the SN2 type; yet the orientation of addition shows that the nucleophile becomes attached to the more highly substituted carbon as though there were a free carbonium ion intermediate. As we shall see, the transition state in reactions of such unstable threemembered rings has much SN1 character. Reduction is generally not stereospecific and can, in certain special cases, be accompanied by rearrangement. Despite the stereospecificity of the first stage, then, the overall process is not,in general, stereospecific. Rearrangements can occur, but are not common. The reaction of 3,3-dimethyl-1-butene illustrates the absence of the rearrangements that are typical of intermediate carbonium ions. Diels-Alder Reaction. The Diels–Alder reaction is a reaction (specifically, a cycloaddition) between a conjugated diene and a substituted alkene, commonly termed the dienophile, to form a substituted cyclohexene system. The reaction can proceed even if some of the atoms in the newly formed ring are not carbon. Some of the Diels–Alder reactions are reversible; the decomposition reaction of the cyclic system is then called the retro-Diels–Alder. The Diels–Alder reaction is generally considered one of the more useful reactions in organic chemistry since it requires very little energy to create a cyclohexene ring, which is useful in many other organic reactions A concerted, single-step mechanism is almost certainly involved; both new carbon-carbon bonds are partly formed in the same transition state, although not necessarily to the same extent. The Diels-Alder reaction is the most important example of cycloaddition. Since reaction involves a system of 4 π electrons (the diene) and a system of 2 π it electrons (the dienophile), it is known as a [4 + 2] cycloaddition. Catalytic addition of hydrogen. Catalytic hydrogenation of alkenes produce the corresponding alkanes. The reaction is carried out under pressure in the presence of a metallic catalyst. Common industrial catalysts are based on platinum, nickel or palladium, but for laboratory syntheses, Raney nickel (formed from an alloy of nickel and aluminium) is often employed. The catalytic hydrogenation of ethylene to yield ethane proceeds thusly: Electrophilic addition. Most addition reactions to alkenes follow the mechanism of electrophilic addition. An example is the Prins reaction, where the electrophile is a carbonyl group. Halogenation. Addition of elementary bromine or chlorine in the presence of an organic solvent to alkenes yield vicinal dibromo- and dichloroalkanes, respectively. The decoloration of a solution of bromine in water is an analytical test for the presence of alkenes: CH2=CH2 + Br2 → BrCH2-CH2Br The reaction works because the high electron density at the double bond causes a temporary shift of electrons in the Br-Br bond causing a temporary induced dipole. This makes the Br closest to the double bond slightly positive and therefore an electrophile. Hydrohalogenation. Addition of hydrohalic acids like HCl or HBr to alkenes yield the corresponding haloalkanes. If the two carbon atoms at the double bond are linked to a different number of hydrogen atoms, the halogen is found preferentially at the carbon with less hydrogen substituents (Markovnikov's rule). Addition of a carbene or carbenoid yields the corresponding cyclopropane Oxidation. Alkenes are oxidized with a large number of oxidizing agents. In the presence of oxygen, alkenes burn with a bright flame to form carbon dioxide and water. Catalytic oxidation with oxygen or the reaction with percarboxylic acids yields epoxides. Reaction with ozone in ozonolysis leads to the breaking of the double bond, yielding two aldehydes or ketones: R1-CH=CH-R2 + O3 → R1-CHO + R2-CHO + H2O This reaction can be used to determine the position of a double bond in an unknown alkene. Polymerization. Polymerization of alkenes is an economically important reaction which yields polymers of high industrial value, such as the plastics polyethylene and polypropylene. Polymerization can either proceed via a free-radical or an ionic mechanism. =Substitution and Elimination Reaction Mechanisms= Nucleophilic Substitution Reactions. Nucleophilic substitution reactions (SN1 and SN2) are very closely related to the E1 and E2 elimination reactions, discussed later in this section, and it is generally a good idea to learn the reactions together, as there are parallels in reaction mechanism, preferred substrates, and the reactions sometimes compete with each other. It's important to understand that substitution and elimination reactions are not associated with a specific compound or mixture so much as they're a representation of how certain reactions take place. At times, combinations of these mechanisms may occur together in the same reaction or may compete against each other, with influences such as solvent or nucleophile choice being the determining factor as to which reaction will dominate. In nucleophilic substitution, a nucleophile attacks a molecule and takes the place of another nucleophile, which then leaves. The nucleophile that leaves is called the leaving group. Nucleophilic substitutions "require " A leaving group is a charged or neutral moiety (group) which breaks free. SN1 vs SN2. One of the main differences between SN1 and SN2 is that the SN1 reaction is a 2-step reaction, initiated by disassociation of the leaving group. The SN2 reaction, on the other hand, is a 1-step reaction where the attacking nucleophile, because of its higher affinity for and stronger bonding with the carbon, forces the leaving group to leave. These two things happen in a single step. These two different mechanisms explain the difference in reaction rates between SN1 and SN2 reactions. SN1 reactions are dependent on the leaving group disassociating itself from the carbon. It is the rate-limiting step and thus, the reaction rate is a first-order reaction whose rate depends solely on that step. Alternatively, in SN2 reactions, the single step of the nucleophile coming together with the reactant from the opposite side of the leaving group, is the key to its rate. Because of this, the rate is dependent on both the concentration of the nucleophile as well as the concentration of the reactant. The higher these two concentrations, the more frequent the collisions. Thus the reaction rate is a second-order reaction: SN2 Reactions. There are primarily 3 things that affect whether an SN2 reaction will take place or not. The most important is structure. That is whether the alkyl halide is on a methyl, primary, secondary, or tertiary carbon. The other two components that determine whether an SN2 reaction will take place or not, are the nucleophilicity of the nucleophile and the solvent used in the reaction. The structure of the alkyl halide has a great effect on mechanism. CH3X &amp; RCH2X are the preferred structures for SN2. R2CHX can undergo the SN2 under the proper conditions (see below), and R3CX rarely, if ever, is involved in SN2 reactions. The reaction takes place by the nucleophile attacking from the opposite side of the bromine atom. Notice that the other 3 bonds are all pointed away from the bromine and towards the attacking nucleophile. When these 3 bonds are hydrogen bonds, there's very little steric hinderance of the approaching nucleophile. However, as the number of R groups increases, so does the steric hinderance, making it more difficult for the nucleophile to get close enough to the α-carbon, to expel the bromine atom. In fact, tertiary carbons (R3CX) are so sterically hindered as to prevent the SN2 mechanism from taking place at all. In the case of this example, a secondary α-carbon, there is still a great deal of steric hinderance and whether the SN2 mechanism will happen will depend entirely on what the nucleophile and solvent are. SN2 reactions are preferred for methyl halides and primary halides. Another important point to keep in mind, and this can be seen clearly in the example above, during an SN2 reaction, the molecule undergoes an inversion. The bonds attached to the α-carbon are pushed away as the nucleophile approaches. During the transition state, these bonds become planar with the carbon and, as the bromine leaves and the nucleophile bonds to the α-carbon, the other bonds fold back away from the nucleophile. This is particularly important in chiral or pro-chiral molecules, where an R configuration will be converted into an S configuration and vice versa. As you'll see below, this is in contrast to the results of SN1 reactions. Examples: OH- is the nucleophile, Cl is the electrophile, HOCH3 is the product, and Cl- is the leaving group. or, The above reaction, taking place in acetone as the solvent, sodium and iodide disassociate almost completely in the acetone, leaving the iodide ions free to attack the CH-Br molecules. The negatively charged iodide ion, a nucleophile, attacks the methyl bromide molecule, forcing off the negatively charged bromide ion and taking its place. The bromide ion is the leaving group. Nucleophilicity. Nucleophilicity is the rate at which a nucleophile displaces the leaving group in a reaction. Generally, nucleophilicity is stronger, the larger, more polarizable, and/or the less stable the nucleophile. No specific number or unit of measure is used. All other things being equal, nucleophiles are generally compared to each other in terms of relative reactivity. For example, a particular strong nucleophile might have a relative reactivity of 10,000 that of a particular weak nucleophile. These relationships are generalities as things like solvent and substrate can affect the relative rates, but they are generally good guidelines for which species make the best nucleophiles. All nucleophiles are Lewis bases. In SN2 reactions, the preferred nucleophile is a strong nucleophile that is a weak base. Examples of these are N3-, RS-, I-, Br-, and CN-. Alternatively, a strong nucleophile that's also a strong base can also work. However, as mentioned earlier in the text, sometimes reaction mechanisms compete and in the case of a strong nucleophile that's a strong base, the SN2 mechanism will compete with the E2 mechanism. Examples of strong nucleophiles that are also strong bases, include RO- and OH-. Leaving Group. Leaving group is the group on the substrate that leaves. In the case of an alkyl halide, this is the halide ion that leaves the carbon atom when the nucleophile attacks. The tendency of the nucleophile to leave is Fluoride ions are very poor leaving groups because they bond very strongly and are very rarely used in alkyl halide substitution reactions. Reactivity of a leaving group is related to its basicity with stronger bases being poorer leaving groups. Solvent. The solvent can play an important role in SN2 reactions, particularly in SN2 involving secondary alkyl halide substrates, where it can be the determining factor in mechanism. Solvent can also have a great effect on reaction rate of SN2 reactions. The SN2 mechanism is preferred when the solvent is an aprotic, polar solvent. That is, a solvent that is polar, but without a polar hydrogen. Polar, protic solvents would include water, alcohols, and generally, solvents with polar NH or OH bonds. Good aprotic, polar solvents are HMPA, CH3CN, DMSO, and DMF. A polar solvent is preferred because it better allows the dissociation of the halide from the alkyl group. A protic solvent with a polar hydrogen, however, forms a 'cage' of hydrogen-bonded solvent around the nucleophile, hindering its approach to the substrate. SN1 Reactions. The SN1 mechanism is very different from the SN2 mechanism. In some of its preferences, its exactly the opposite and, in some cases, the results of the reaction can be significantly different. Like the SN2 mechanism, structure plays an important role in the SN1 mechanism. The role of structure in the SN1 mechanism, however, is quite different and because of this, the reactivity of structures is more or less reversed. The SN1 mechanism is preferred for tertiary alkyl halides and, depending on the solvent, may be preferred in secondary alkyl halides. The SN1 mechanism does not operate on primary alkyl halides or methyl halides. To understand why this is so, let's take a look at how the SN1 mechanism works. At the top of the diagram, the first step is the spontaneous dissociation of the halide from the alkyl halide. Unlike the SN2 mechanism, where the attacking nucleophile causes the halide to leave, the SN1 mechanism depends on the ability of the halide to leave on its own. This requires certain conditions. In particular, the stability of the carbocation is crucial to the ability of the halide to leave. Since we know tertiary carbocations are the most stable, they are the best candidates for the SN1 mechanism. And with appropriate conditions, secondary carbocations will also operate by the SN1 mechanism. Primary and methyl carbocations however, are not stable enough to allow this mechanism to happen. Once the halide has dissociated, the water acts as a nucleophile to bond to the carbocation. In theSN2 reactions, there is an inversion caused by the nucleophile attacking from the opposite side while the halide is still bonded to the carbon. In the SN1 mechanism, since the halide has left, and the bonds off of the α-carbon have become planar, the water molecule is free to attack from either side. This results in, primarily, a racemic mixture. In the final step, one of the hydrogens of the bonded water molecule is attacked by another water molecule, leaving an alcohol. "Note: Racemic mixtures imply entirely equal amounts of mixture, however this is rarely the case in SN1. There is a slight tendency towards attack from the opposite side of the halide. This is the result some steric hinderence from the leaving halide which is sometimes close enough to the leaving side to block the nucleophile's approach from that side." Solvent. Like the SN2 mechanism, the SN1 is affected by solvent as well. As with structure, however, the reasons differ. In the SN1 mechanism, a polar, protic solvent is used. The polarity of the solvent is associated with the dielectric constant of the solvent and solutions with high dielectric constants are better able to support separated ions in solution. In SN2 reactions, we were concerned about polar hydrogen atoms "caging" our nucleophile. This still happens with a polar protic solvent in SN1 reactions, so why don't we worry about it? You have to keep in mind the mechanism of the reaction. The first step, and more importantly, the rate-limiting step, of the SN1 reaction, is the ability to create a stable carbocation by getting the halide anion to leave. With a polar protic solvent, just as with a polar aprotic solvent,we're creating a stable cation, however it's the polar hydrogens that stabilize the halide anion and make it better able to leave. Improving the rate-limiting step is always the goal. The "caging" of the nucleophile is unrelated to the rate-limiting step and even in its "caged" state, the second step, the attack of the nucleophile, is so much faster than the first step, that the "caging" can simply be ignored. Summary. SN1, SN2, E1, and E2, are all reaction mechanisms, not reactions themselves. They are mechanisms used by a number of different reactions. Usually in organic chemistry, the goal is to synthesize a product. In cases where you have possibly competing mechanisms, and this is particularly the case where an SN1 and an E1 reaction are competing, the dominating mechanism is going to decide what your product is, so knowing the mechanisms and which conditions favor one over the other, will determine your product. In other cases, knowing the mechanism allows you to set up an environment favorable to that mechanism. This can mean the difference between having your product in a few minutes, or sometime around the next ice age. So when you're designing a synthesis for a product, you need to consider, I want to get product Y, so what are my options to get to Y? Once you know your options and you've decided on a reaction, then you need to consider the mechanism of the reaction and ask yourself, how do I create conditions that are going to make this happen correctly and happen quickly? Elimination Reactions. Nucleophilic substitution reactions and Elimination reactions share a lot of common characteristics, on top of which, the E1 and SN1 as well as E2 and SN2 reactions can sometimes compete and, since their products are different, it's important to understand them both. Without understanding both kinds of mechanisms, it would be difficult to get the product you desire from a reaction. In addition, the SN1 and SN2 reactions will be referenced quite a bit by way of comparison and contrast, so it's probably best to read that section first and then continue here. Elimination reactions are the mechanisms for creating alkene products from haloalkane reactants. E1 and E2 elimination, unlike SN1 and SN2 substitution, mechanisms do not occur with methyl halides because the reaction creates a double bond between two carbon atoms and methylhalides have only one carbon. E1 vs E2. Reaction rates. E1 and E2 are two different pathways to creating alkenes from haloalkanes. As with SN1 and SN2 reactions, one of the key differences is in the reaction rate, as it provides great insight into the mechanisms. E1 reactions, like SN1 reactions are 2-step reactions. Also like SN1 reactions, the rate-limiting step is the dissociation of the halide from its alkane, making it a first-order reaction, depending on the concentration of the haloalkane, with a reaction rate of: On the other hand, E2 reactions, like SN2 reactions are 1-step reactions. And again, as with SN2 reactions, the rate limiting step is the ability of a nucleophile to attach to the alkane and displace the halide. Thus it is a second-order reaction that depends on the concentrations of both the nucleophile and haloalkane, with a reaction rate of: Zaitsev's Rule. Zaitsev's rule (sometimes spelled "Saytzeff") states that in an elimination reaction, when multiple products are possible, the most stable alkene is the major product. That is to say, the most highly substituted alkene (the alkene with the most non-hydrogen substituents) is the major product. Both E1 and E2 reactions produce a mixture of products, when possible, but generally follow Zaitsev's rule. We'll see below why E1 reactions follow Zaitsev's rule more reliably and tend to produce a purer product. The above image represents two possible pathways for the dehydrohalogenation of (S)-2-bromo-3-methylbutane. The two potential products are 2-methylbut-2-ene and 3-methylbut-1-ene. The images on the right are simplified drawings of the molecular product shown in the images in the center. As you can see on the left, the bromine is on the second carbon and in an E1 or E2 reaction, the hydrogen could be removed from either the 1st or the 3rd carbon. Zaitsev's rule says that the hydrogen will be removed predominantly from the 3rd carbon. In reality, there will be a mixture, but most of the product will be 2-methylbut-2-ene by the E1 mechanism. By the E2 reaction, as we'll see later, this might not necessarily be the case. E2. The E2 mechanism is concerted and highly stereospecific, because it can occur only when the H and the leaving group X are in an anti-coplanar position. That is, in a Newman projection, the H and X must be 180°, or in the anti-configuration. This behaviour stems from the best overlap of the 2"p" orbitals of the adjacent carbons when the pi bond has to be formed. If the H and the leaving group cannot be brought into this position due to the structure of the molecule, the "E2" mechanism will not take place. Therefore, only molecules having accessible H-X anti-coplanar conformations can react via this route. Furthermore, the E2 mechanism will operate contrary to Zaitsev's rule if the only anti-coplanar hydrogen from the leaving group results in the least stable alkene. A good example of how this can happen is by looking at how cyclohexane and cyclohexene derivatives might operate in E2 conditions. Let's look at the example above. The reactant we're using is 1-chloro-2-isopropylcyclohexane. The drawing at the top left is one conformation and the drawing below is after a ring flip. In the center are Newman projections of both conformations and the drawings on the right, the products. If we assume we're treating the 1-chloro-2-isopropylcyclohexane with a strong base, for example CH3CH2O- (ethanolate), the mechanism that dominates is E2. There are 3 hydrogens off of the carbons adjacent to our chlorinated carbon. The red and the green ones are two of them. The third would be hard to show but is attached to the same carbon as the red hydrogen, angled a little down from the plane and towards the viewer. The red hydrogen is the only hydrogen that's 180° from the chlorine atom, so it's the only one eligible for the E2 mechanism. Because of this, the product is going to be only 3-isopropylcylcohex-1-ene. Notice how this is contrary to Zaitsev's rule which says the most substituted alkene is preferred. By his rule, 1-isopropylcyclohexene should be our primary product, as that would leave the most substituted alkene. However it simply can't be produced because of the steric hindrance. The images below shows the molecule after a ring flip. In this conformation, no product is possible. As you can see from the Newman projection, there are no hydrogens 180° from the chlorine atom. So it's important, when considering the E2 mechanism, to understand the geometry of the molecule. Sometimes the geometry can be used to your advantage to preferentially get a single product. Other times it will prevent you from getting the product you want, and you'll need to consider a different mechanism to get your product. "Note: Often the word periplanar is used instead of coplanar. Coplanar implies precisely 180 degree separation and "peri-", from Greek for "near", implies near 180 degrees. Periplanar may actually be more accurate. In the case of the 1-chloro-3-isopropylcyclohexane example, because of molecular forces, the chlorine atom is actually slightly less than 180 degrees from both the hydrogen and the isopropyl group, so in this case, periplanar might be a more correct term." E1. The E1 mechanism begins with the dissociation of the leaving group from an alkyl, producing a carbocation on the alkyl group and a leaving anion. This is the same way the SN1 reaction begins, so the same thing that helps initiate that step in SN1 reactions, help initiate the step in E1 reactions. More specifically, secondary and tertiary carbocations are preferred because they're more stable than primary carbocations. The choice of solvent is the same as SN1 as well; a polar protic solvent is preferred because the polar aspect stabilizes the carbocation and the protic aspect stabilizes the anion. What makes the difference between whether the reaction takes the SN1 or E1 pathway then, must depend on the second step; the action of the nucleophile. In SN1 reactions, a strong nucleophile that's a weak base is preferred. The nucleophile will then attack and bond to the carbocation. In E1 reactions, a strong nucleophile is still preferred. The difference is that a strong nucleophile that's also a strong base, causes the nucleophile to attack the hydrogen at the β-carbon instead of the α-carbocation. The nucleophile/base then extracts the hydrogen causing the bonding electrons to fall in and produce a pi bond with the carbocation. Because the hydrogen and the leaving group are lost in two separate steps and the fact that it has no requirements as to geometry, the E1 mechanism more reliably produces products that follow Zaitsev's rule. =References= « Haloalkanes |Alkenes| Alkynes » 

Introduction. Chirality (pronounced kie-RAL-it-tee) is the property of "handedness". If you attempt to superimpose your right hand on top of your left, the two do not match up in the sense that your right hand's thumb overlays your left hand's pinky finger. Your two hands cannot be superimposed identically, despite the fact that your fingers of each hand are connected in the same way. Any object can have this property, including molecules. An object that is chiral is an object that can not be superimposed on its mirror image. Chiral objects don't have a "plane of symmetry". An achiral object has a plane of symmetry or a rotation-reflection axis, i.e. reflection gives a rotated version. Optical isomers or enantiomers are stereoisomers which exhibit chirality. Optical isomerism is of interest because of its application in inorganic chemistry, organic chemistry, physical chemistry, pharmacology and biochemistry. They are often formed when asymmetric centers are present, for example, a carbon with four different groups bonded to it. Every stereocenter in one enantiomer has the opposite configuration in the other. When a molecule has more than one source of asymmetry, two optical isomers may be neither perfect reflections of each other nor superimposeable: some but not all stereocenters are inverted. These molecules are an example of diastereomers: They are not enantiomers. Diastereomers seldom have the same physical properties. Sometimes, the stereocentres are themselves symmetrical. This causes the counterintuitive situation where two chiral centres may be present but no isomers result. Such compounds are called meso compounds. A mixture of equal amounts of both enantiomers is said to be a racemic mixture. It is the symmetry of a molecule (or any other object) that determines whether it is chiral or not. Technically, a molecule is achiral (not chiral) if and only if it has an axis of improper rotation; that is, an n-fold rotation (rotation by 360°/n) followed by a reflection in the plane perpendicular to this axis which maps the molecule onto itself. A chiral molecule is not necessarily dissymmetric (completely devoid of symmetry) as it can have, e.g., rotational symmetry. A simplified rule applies to tetrahedrally-bonded carbon, as shown in the illustration: if all four substituents are different, the molecule is chiral. It is important to keep in mind that molecules which are dissolved in solution or are in the gas phase usually have considerable flexibility and thus may adopt a variety of different conformations. These various conformations are themselves almost always chiral. However, when assessing chirality, one must use a structural picture of the molecule which corresponds to just one chemical conformation – the one of lowest energy. Chiral Compounds With Stereocenters. Most commonly, chiral molecules have point chirality, centering around a single atom, usually carbon, which has four different substituents. The two enantiomers of such compounds are said to have different absolute configurations at this center. This center is thus stereogenic (i.e., a grouping within a molecular entity that may be considered a focus of stereoisomerism), and is exemplified by the α-carbon of amino acids. The special nature of carbon, its ability to form four bonds to different substituents, means that a mirror image of the carbon with four different bonds will not be the same as the original compound, no matter how you try to rotate it. Understanding this is vital because the goal of organic chemistry is understanding how to use tools to synthesize a compound with the desired chirality, because a different arrangement may have no effect, or even an undesired one. A carbon atom is chiral if it has four different items bonded to it at the same time. Most often this refers to a carbon with three heteroatoms and a hydrogen, or two heteroatoms plus a bond to another carbon plus a bond to a hydrogen atom. It can also refer to a nitrogen atom bonded to four different types of molecules, if the nitrogen atom is utilizing its lone pair as a nucleophile. If the nitrogen has only three bonds it is not chiral, because the lone pair of electrons can flip from one side of the atom to the other spontaneously. "Any atom in an organic molecule that is bonded to four different types of atoms or chains of atoms can be considered "chiral"." If a carbon atom (or other type of atom) has four different substituents, that carbon atom forms a "chiral center" (also known as a "stereocenter"). Chiral molecules often have one or more stereocenters. When drawing molecules, stereocenters are usually indicated with an asterisk near the carbon. Example: Left: The carbon atom has a Cl, a Br, and 2 CH3. That's only 3 different substituents, which means this is not a stereocenter. Center: The carbon atom has one ethyl group (CH2CH3), one methyl group (CH3) and 2 H. This is not a stereocenter. Right: The carbon atom has a Cl and 1 H. Then you must look around the ring. Since one side has a double bond and the other doesn't, it means the substituents off that carbon are different. The 4 different substituents make this carbon a stereocenter and makes the molecule chiral. A molecule can have multiple chiral centers without being chiral overall: It is then called a meso compound. This occurs if there is a symmetry element (a mirror plane or inversion center) which relates the chiral centers. Fischer projections. Fischer projections (after the German chemist ) are an ingenious means for representing configurations of carbon atoms. Considering the carbon atom as the center, the bonds which extend towards the viewer are placed horizontally. Those extending away from the viewer are drawn vertically. This process, when using the common dash and wedge representations of bonds, yields what is sometimes referred to as the "bowtie" drawing due to its characteristic shape. This representation is then further shorthanded as two lines: the horizontal (forward) and the vertical (back), as showed in the figure below: Naming conventions. There are three main systems for describing configuration: the oldest, the "relative" whose use is now deprecated, and the current, or "absolute". The relative configuration description is still used mainly in glycochemistry. Configuration can also be assigned on the purely empirical basis of the optical activity. By optical activity: (+)- and (-)-. An optical isomer can be named by the direction in which it rotates the plane of polarized light. If an isomer rotates the plane clockwise as seen by a viewer towards whom the light is traveling, that isomer is labeled (+). Its counterpart is labeled (-). The (+) and (-) isomers have also been termed d- and l-, respectively (for dextrorotatory and levorotatory). This labeling is easy to confuse with D- and L- and is therefore not encouraged by IUPAC. The fact that an enantiomer can rotate polarised light clockwise ("d"- or "+"- enantiomer) does not relate with the relative configuration (D- or L-) of it. By relative configuration: D- and L-. Fischer, whose research interest was in carbohydrate chemistry, took glyceraldehyde (the simplest sugar, systematic name 2,3-dihydroxypropanal) as a template chiral molecule and denoted the two possible configurations with D- and L-, which rotated polarised light clockwise and counterclockwise, respectively. All other molecules are assigned the D- or L- configuration if the chiral centre can be formally obtained from glyceraldehyde by substitution. For this reason the D- or L- naming scheme is called "relative configuration". An optical isomer can be named by the spatial configuration of its atoms. The D/L system does this by relating the molecule to glyceraldehyde. Glyceraldehyde is chiral itself, and its two isomers are labeled D and L. Certain chemical manipulations can be performed on glyceraldehyde without affecting its configuration, and its historical use for this purpose (possibly combined with its convenience as one of the smallest commonly-used chiral molecules) has resulted in its use for nomenclature. In this system, compounds are named by analogy to glyceraldehyde, which generally produces unambiguous designations, but is easiest to see in the small biomolecules similar to glyceraldehyde. One example is the amino acid alanine: alanine has two optical isomers, and they are labeled according to which isomer of glyceraldehyde they come from. Glycine, the amino acid derived from glyceraldehyde, incidentally, does not retain its optical activity, since its central carbon is not chiral. Alanine, however, is essentially methylated glycine and shows optical activity. The D/L labeling is unrelated to (+)/(-); it does not indicate which enantiomer is dextrorotatory and which is levorotatory. Rather, it says that the compound's stereochemistry is related to that of the dextrorotatory or levorotatory enantiomer of glyceraldehyde. Nine of the nineteen L-amino acids commonly found in proteins are dextrorotatory (at a wavelength of 589 nm), and D-fructose is also referred to as levulose because it is levorotatory. The dextrorotatory isomer of glyceraldehyde is in fact the D isomer, but this was a lucky guess. At the time this system was established, there was no way to tell which configuration was dextrorotatory. (If the guess had turned out wrong, the labeling situation would now be even more confusing.) A rule of thumb for determining the D/L isomeric form of an amino acid is the "CORN" rule. The groups: are arranged around the chiral center carbon atom. If these groups are arranged clockwise around the carbon atom, then it is the L-form. If counter-clockwise, it is the D-form.This rule only holds when the hydrogen atom is pointing out of the page. By absolute configuration: R- and S-. Main article: R-S System The absolute configuration system stems from the , which allow a precise description of a stereocenter without using any reference compound. In fact the basis is now the atomic number of the stereocenter substituents. The R/S system is another way to name an optical isomer by its configuration, without involving a reference molecule such as glyceraldehyde. It labels each chiral center R or S according to a system by which its ligands are each assigned a priority, according to the Cahn Ingold Prelog priority rules, based on atomic number. This system labels each chiral center in a molecule (and also has an extension to chiral molecules not involving chiral centers). It thus has greater generality than the D/L system, and can label, for example, an (R,R) isomer versus an (R,S) — diastereomers. The R/S system has no fixed relation to the (+)/(-) system. An R isomer can be either dextrorotatory or levorotatory, depending on its exact ligands. The R/S system also has no fixed relation to the D/L system. For example, one of glyceraldehyde's ligands is a hydroxy group, -OH. If a thiol group, -SH, were swapped in for it, the D/L labeling would, by its definition, not be affected by the substitution. But this substitution would invert the molecule's R/S labeling, due to the fact that sulfur's atomic number is higher than carbon's, whereas oxygen's is lower. [Note: This seems incorrect. Oxygen has a higher atomic number than carbon. Sulfur has a higher atomic number than oxygen. The reason the assignment priorities change in this example is because the CH2SH group gets a higher priority than the CHO, whereas in glyceraldehyde the CHO takes priority over the CH2OH.] For this reason, the D/L system remains in common use in certain areas, such as amino acid and carbohydrate chemistry. It is convenient to have all of the common amino acids of higher organisms labeled the same way. In D/L, they are all L. In R/S, they are not, conversely, all S — most are, but cysteine, for example, is R, again because of sulfur's higher atomic number. The word “racemic” is derived from the Latin word for grape; the term having its origins in the work of Louis Pasteur who isolated racemic tartaric acid from wine. Chiral Compounds Without Stereocenters. It is also possible for a molecule to be chiral without having actual point chirality (stereocenters). Commonly encountered examples include 1,1'-bi-2-naphthol (BINOL) and 1,3-dichloro-allene which have axial chirality, and (E)-cyclooctene which has planar chirality. For example, the isomers which are shown by the following figure are different. The two isomers cannot convert from one to another spontaneously because of restriction of rotation of double bonds. Other types of chiral compounds without stereocenters (like restriction of rotation of a single bond because of steric hindrance) also exist. Consider the following example of the R and S binol molecules: "The biphenyl C-C bond cannot rotate if the X and Y groups cause steric hindrance." "This compound exhibits spiral chirality." Properties of optical isomers. Enantiomers have – "when present in a symmetric environment" – identical chemical and physical properties except for their ability to rotate plane-polarized light by equal amounts but in opposite directions. A solution of equal parts of an optically-active isomer and its enantiomer is known as a racemic solution and has a net rotation of plane-polarized light of zero. Enantiomers differ in how they interact with different optical isomers of other compounds. In nature, most biological compounds (such as amino acids) occur as single enantiomers. As a result, different enantiomers of a compound may have substantially different biological effects. Different enantiomers of the same chiral drug can have very different pharmological effects, mainly because the proteins they bind to are also chiral. For example, spearmint leaves and caraway seeds respectively contain L-carvone and D-carvone – enantiomers of carvone. These smell different to most people because our taste receptors also contain chiral molecules which behave differently in the presence of different enantiomers. D-form Amino acids tend to taste sweet, whereas L-forms are usually tasteless. This is again due to our chiral taste molecules. The smells of oranges and lemons are examples of the D and L enantiomers. Penicillin's activity is stereoselective. The antibiotic only works on peptide links of D-alanine which occur in the cell walls of bacteria – but not in humans. The antibiotic can kill only the bacteria, and not us, because we don't have these D-amino acids. The electric and magnetic fields of polarized light oscillate in a geometric plane. An axis normal to this plane gives the direction of energy propagation. Optically active isomers rotate the plane that the fields oscillate in. The polarized light is actually rotated in a racemic mixture as well, but it is rotated to the left by one of the two enantiomers, and to the right by the other, which cancel out to zero net rotation. Chirality in biology. Many biologically-active molecules are chiral, including the naturally-occurring amino acids (the building blocks of proteins), and sugars. Interestingly, in biological systems most of these compounds are of the same chirality: most amino acids are L and sugars are D. The origin of this homochirality in biology is the subject of much debate. Chiral objects have different interactions with the two enantiomers of other chiral objects. Enzymes, which are chiral, often distinguish between the two enantiomers of a chiral substrate. Imagine an enzyme as having a glove-like cavity which binds a substrate. If this glove is right handed, then one enantiomer will fit inside and be bound while the other enantiomer will have a poor fit and is unlikely to bind. Chirality in inorganic chemistry. Many coordination compounds are chiral; for example the well-known [Ru(2,2'-bipyridine)3]2+ complex in which the three bipyridine ligands adopt a chiral propeller-like arrangement [7]. In this case, the Ru atom may be regarded as a stereogenic centre, with the complex having point chirality. The two enantiomers of complexes such as [Ru(2,2'-bipyridine)3]2+ may be designated as Λ (left-handed twist of the propeller described by the ligands) and Δ (right-handed twist). Hexol is a chiral cobalt compound. Enantiopure preparations. Several strategies exist for the preparation of enantiopure compounds. The first method is the separation of a racemic mixture into its isomers. Louis Pasteur in his pioneering work was able to isolate the isomers of tartaric acid because they crystallize from solution as crystals with differing symmetry. A less common and more recently discovered method is by enantiomer self-disproportionation, which is an advanced technique involving the separation of a primarily racemic fraction from a nearly enantiopure fraction via column chromatography. In a non-symmetric environment (such as a biological environment) enantiomers may react at different speeds with other substances. This is the basis for "chiral synthesis", which preserves a molecule's desired chirality by reacting it with or catalyzing it with chiral molecules capable of maintaining the product's chirality in the desired conformation (using certain chiral molecules to help it keep its configuration). Other methods also exist and are used by organic chemists to synthesize only (or maybe only "mostly") the desired enantiomer in a given reaction. Enantiopure medications. Advances in industrial chemical processes have allowed pharmaceutical manufacturers to take drugs that were originally marketed in racemic form and divide them into individual enantiomers, each of which may have unique properties. For some drugs, such as zopiclone, only one enantiomer (eszopiclone) is active; the FDA has allowed such once-generic drugs to be patented and marketed under another name. In other cases, such as ibuprofen, both enantiomers produce the same effects. Steroid receptor sites also show stereoisomer specificity. Examples of racemic mixtures and enantiomers that have been marketed include: Many chiral drugs must be made with high enantiomeric purity due to potential side-effects of the other enantiomer. (The other enantiomer may also merely be inactive.) Consider a racemic sample of thalidomide. One enantiomer was thought to be effective against morning sickness while the other is now known to be teratogenic. Unfortunately, in this case administering just one of the enantiomers to a pregnant patient would still be very dangerous as the two enantiomers are readily interconverted "in vivo". Thus, if a person is given either enantiomer, both the D and L isomers will eventually be present in the patient's serum and so chemical processes may not be used to mitigate its toxicity. See also. Optical activity 

In alkene chemistry, we demonstrated that allylic carbon could maintain a cation charge because the double bond could de-localize to support the charge. What of having two double bonds separated by a single bond? What of having a compound that alternates between double bond and single bond? In addition to other concepts, this chapter will explore what a having a conjugated system means in terms of stability and reaction. Dienes are simply hydrocarbons which contain two double bonds. Dienes are intermediate between alkenes and polyenes. Dienes can divided into three classes: 

=Resources= =Other online textbooks= 

Brief History. Jöns Jacob Berzelius, a physician by trade, first coined the term "organic chemistry" in 1806 for the study of compounds derived from biological sources. Up through the early 19th century, naturalists and scientists observed critical differences between compounds that were derived from living things and those that were not. Chemists of the period noted that there seemed to be an essential yet inexplicable difference between the properties of the two different types of compounds. The vital force theory, sometimes called "vitalism" (vital means "life force"), was therefore proposed, and widely accepted, as a way to explain these differences, that a "vital force" existed within organic material but did not exist in any inorganic materials. Synthesis of Urea. Friedrich Wöhler is widely regarded as a pioneer in organic chemistry as a result of his synthesizing of the biological compound urea (a component of urine in many animals) utilizing what is now called "the Wöhler synthesis." Wöhler mixed silver or lead cyanate with ammonium nitrate; this was supposed to yield ammonium cyanate as a result of an exchange reaction, according to Berzelius's dualism theory. Wöhler, however, discovered that the end product of this reaction is "not" ammonium cyanate (NH4OCN), an inorganic salt, but urea ((NH2)2CO), a biological compound. (Furthermore, heating ammonium cyanate turns it into urea.) Faced with this result, Berzelius had to concede that (NH2)2CO and NH4OCN were "isomers". Until this discovery in the year 1828, it was widely believed by chemists that organic substances could only be formed under the influence of the "vital force" in the bodies of animals and plants. Wöhler's synthesis dramatically proved that view to be false. Urea synthesis was a critical discovery for biochemists because it showed that a compound known to be produced in nature only by biological organisms could be produced in a laboratory under controlled conditions from inanimate matter. This "in vitro" synthesis of organic matter disproved the common theory (vitalism) about the vis vitalis, a transcendent "life force" needed for producing organic compounds. Organic vs Inorganic Chemistry. Although originally defined as the chemistry of biological molecules, organic chemistry has since been redefined to refer specifically to carbon compounds — even those with non-biological origin. Some carbon molecules are not considered organic, with carbon dioxide being the most well known and most common inorganic carbon compound, but such molecules are the exception and not the rule. Organic chemistry focuses on carbon and following movement of the electrons in carbon chains and rings, and also how electrons are shared with other carbon atoms and heteroatoms. Organic chemistry is primarily concerned with the properties of covalent bonds and non-metallic elements, though ions and metals do play critical roles in some reactions. The applications of organic chemistry are myriad, and include all sorts of plastics, dyes, flavorings, scents, detergents, explosives, fuels and many, many other products. Read the ingredient list for almost any kind of food that you eat — or even your shampoo bottle — and you will see the handiwork of organic chemists listed there. Major Advances in the Field of Organic Chemistry. Of course a chemistry text should at least mention Antoine Laurent Lavoisier. The French chemist is often called the "Father of Modern Chemistry" and his place is first in any pantheon of great chemistry figures. Your general chemistry textbook should contain information on the specific work and discoveries of Lavoisier — they will not be repeated here because his discoveries did not relate directly to organic chemistry in particular. Berzelius and Wöhler are discussed above, and their work was foundational to the specific field of organic chemistry. After those two, three more scientists are famed for independently proposing the elements of structural theory. Those chemists were August Kekulé, Archibald Couper, and Alexander Butlerov. Kekulé was a German, an architect by training, and he was perhaps the first to propose that isomerism was due to carbon's proclivity towards forming four bonds. Its ability to bond with up to four other atoms made it ideal for forming long chains of atoms in a single molecule, and also made it possible for the same number of atoms to be connected in an enormous variety of ways. Couper, a Scot, and Butlerov, a Russian, came to many of the same conclusions at the same time or just a short time after. Through the nineteenth century and into the twentieth, experimental results brought to light much new knowledge about atoms, molecules, and molecular bonding. In 1916 it was Gilbert Lewis of U.C. Berkeley who described covalent bonding largely as we know it today (electron-sharing). Nobel laureate Linus Pauling further developed Lewis' concepts by proposing resonance while he was at the California Institute of Technology. At about the same time, Sir Robert Robinson of Oxford University focused primarily on the electrons of atoms as the engines of molecular change. Sir Christopher Ingold of University College, London, organized what was known of organic chemical reactions by arranging them in schemes we now know as mechanisms, in order to better understand the sequence of changes in a synthesis or reaction. The field of organic chemistry is probably the most active and important field of chemistry at the moment, due to its extreme applicability to both biochemistry (especially in the pharmaceutical industry) and petrochemistry (especially in the energy industry). Organic chemistry has a relatively recent history, but it will have an enormously important future, affecting the lives of everyone around the world for many, many years to come. &lt;noinclude&gt; « Foundational concepts| History of Organic Chemistry | Atomic Structure &gt; | Alkanes » &lt;/alkynes&gt; 

List of Authors. "You always need someone to blame." 

Atomic Structure. Atoms are made up of a nucleus and electrons that orbit the nucleus. The nucleus consists of protons and neutrons. An atom in its natural, uncharged state has the same number of electrons as protons. The nucleus. The nucleus is made up of protons, which are positively charged, and neutrons, which have no charge. Neutrons and protons have about the same mass, and together account for most of the mass of the atom. Electrons. The electrons are negatively charged particles. The mass of an electron is about 2000 times smaller than that of a proton or neutron at 0.00055 amu. Electrons circle so fast that it cannot be determined where electrons are at any point in time. The image depicts the old Bohr model of the atom, in which the electrons inhabit discrete "orbitals" around the nucleus much like planets orbit the sun. This model is outdated. Current models of the atomic structure hold that electrons occupy fuzzy clouds around the nucleus of specific shapes, some spherical, some dumbbell shaped, some with even more complex shapes. Shells and Orbitals. Electron shells. Electrons orbit atoms in clouds of distinct shapes and sizes. The electron clouds are layered one inside the other into units called shells, with the electrons occupying the simplest orbitals in the innermost shell having the lowest energy state and the electrons in the most complex orbitals in the outermost shell having the highest energy state. The higher the energy state, the more energy the electron has, just like a rock at the top of a hill has more potential energy than a rock at the bottom of a valley. The main reason why electrons exist in higher energy orbitals is because only two electrons can exist in any orbital. So electrons fill up orbitals, always taking the lowest energy orbitals available. An electron can also be pushed to a higher energy orbital, for example by a photon. Typically this is not a stable state and after a while the electron descends to lower energy states by emitting a photon spontaneously. These concepts will be important in understanding later concepts like optical activity of chiral compounds as well as many interesting phenomena outside the realm of organic chemistry (for example, how lasers work). Electron orbitals. Each different shell is subdivided into one or more orbitals, which also have different energy levels, although the energy difference between orbitals is less than the energy difference between shells. Longer wavelengths have less energy; the s orbital has the longest wavelength allowed for an electron orbiting a nucleus and this orbital is observed to have the lowest energy. Each orbital has a characteristic shape which shows where electrons most often exist. The orbitals are named using letters of the alphabet. In order of increasing energy the orbitals are: s, p, d, and f orbitals. As one progresses up through the shells (represented by the principal quantum number n) more types of orbitals become possible. The shells are designated by numbers. So the 2s orbital refers to the s orbital in the second shell. S orbital. The s orbital is the orbital lowest in energy and is spherical in shape. Electrons in this orbital are in their fundamental frequency. This orbital can hold a maximum of two electrons. P orbital. The next lowest-energy orbital is the p orbital. Its shape is often described as like that of a dumbbell. There are three p-orbitals each oriented along one of the 3-dimensional coordinates x, y or z. Each of these three "p" orbitals can hold a maximum of two electrons. These three different p orbitals can be referred to as the px, py, and pz. The s and p orbitals are important for understanding most of organic chemistry as these are the orbitals that are occupied in the type of atoms that are most common in organic compounds. D and F orbitals. There are also D and F orbitals. D orbitals are present in transition metals. Sulphur and phosphorus have empty D orbitals. Compounds involving atoms with D orbitals do come into play, but are rarely part of an organic molecule. F are present in the elements of the lanthanide and actinide series. Lanthanides and actinides are mostly irrelevant to organic chemistry. Filling electron shells. When an atom or ion receives electrons into its orbitals, the orbitals and shells fill up in a particular manner. There are three principles that govern this process: Pauli exclusion principle. No two electrons in an atom can have all four quantum numbers the same. What this translates to in terms of our pictures of orbitals is that each orbital can only hold two electrons, one "spin up" and one "spin down". Hund's rule. This states that filled and half-filled shells tend to have additional stability. In some instances, then, for example, the 4s orbitals will be filled before the 3d orbitals. This rule is applicable only for those elements that have d electrons, and so is less important in organic chemistry (though it is important in organometallic chemistry). Octet rule. The octet rule states that atoms tend to prefer to have eight electrons in their valence shell, so will tend to "combine" in such a way that each atom can have eight electrons in its valence shell, similar to the electronic configuration of a noble gas. In simple terms, molecules are more stable when the outer shells of their constituent atoms are empty, full, or have eight electrons in the outer shell. The main exception to the rule is helium, which is at lowest energy when it has two electrons in its valence shell. Other notable exceptions are aluminum and boron, which can function well with six valence electrons; and some atoms beyond group three on the periodic table that can have over eight electrons, such as sulphur. Additionally, some noble gasses can form compounds when expanding their valence shell. Hybridization. Hybridization refers to the combining of the orbitals of two or more covalently bonded atoms. Depending on how many free electrons a given atom has and how many bonds it is forming, the electrons in the s and the p orbitals will combine in certain manners to form the bonds. It is easy to determine the hybridization of an atom given a Lewis structure. First, you count the number of pairs of free electrons and the number of "sigma" bonds (single bonds). Do not count double bonds, since they do not affect the hybridization of the atom. Once the total of these two is determined, the hybridization pattern is as follows:  Sigma Bonds + Electron Pairs Hybridization  2 sp  3 sp2  4 sp3 The pattern here is the same as that for the electron orbitals, which serves as a memory guide. 

Ionic Bonding. Ionic bonding is when positively and negatively charged ions stick to each other through electrostatic force. These bonds are " slightly weaker than covalent bonds" and stronger than Van der Waals bonding or hydrogen bonding. In ionic bonds the electronegativity of the negative ion is so much stronger than the electronegativity of the positive ion that the two ions do not share electrons. Rather, the more electronegative ion assumes full ownership of the electron(s). Perhaps the most common example of an ionically bonded substance is NaCl, or table salt. In this, the sodium (Na) atom gives up an electron to the much more electronegative chlorine (Cl) atom, and the two atoms become ions, Na+ and Cl-.The electrostatic bonding force between the two oppositely charged ions extends outside the local area attracting other ions to form giant crystal structures. For this reason most ionically bonded materials are solid at room temperature. Sodium chloride forms crystals with cubic symmetry. In these, the larger chloride ions are arranged in a cubic close-packing, while the smaller sodium ions fill the octahedral gaps between them. Each ion is surrounded by six of the other kind. This same basic structure is found in many other minerals, and is known as the halite structure. Covalent Bonding. Covalent bonding is close to the heart of organic chemistry. This is where two atoms share electrons in a bond. The goal of each atom is to "fill its octet" as well as have a "formal charge of zero". To do this, atomic nuclei share electrons in the space between them. This sharing also allows the atoms to reach a lower energy state, which stabilizes the molecule. Most reactions in chemistry are due to molecules achieving a lower energy state. Covalent bonds are most frequently seen between atoms with similar electronegativity. In molecules that only have one type of atom, e.g. H2 or O2 , the electronegativity of the atoms is essentially identical, so they cannot form ionic bonds. They always form covalent bonds. Carbon is especially good at covalent bonding because its electronegativity is intermediate relative to other atoms. That means it can give as well as take electrons as needs warrant. Covalently bonded compounds have strong internal bonds but weak attractive forces between molecules. Because of these weak attractive forces, the melting and boiling points of these compounds are much lower than compounds with ionic bonds. Therefore, such compounds are much more likely to be liquids or gases at room temperature than ionically bonded compounds. In molecules formed from two atoms of the same element, there is no difference in the electronegativity of the bonded atoms, so the electrons in the covalent bond are shared equally, resulting in a completely non-polar covalent bond. In covalent bonds where the bonded atoms are different elements, there is a difference in electronegativities between the two atoms. The atom that is more electronegative will attract the bonding electrons more toward itself than the less electronegative atom. The difference in charge on the two atoms because of the electrons causes the covalent bond to be polar. Greater differences in electronegativity result in more polar bonds. Depending on the difference in electronegativities, the polarity of a bond can range from non-polar covalent to ionic with varying degrees of polar covalent in between. An overall imbalance in charge from one side of a molecule to the other side is called a dipole moment. Such molecules are said to be polar. For a completely symmetrical covalently bonded molecule, the overall dipole moment of the molecule is zero. Molecules with larger dipole moments are more polar. The most common polar molecule is water. Bond Polarity and Dipole Moment. The ideas of bond polarity and dipole moment play important roles in organic chemistry. If you look at the image of methane on the right, the single most important aspect of it in terms of bond polarity is that it is a symmetric molecule. It has 4 hydrogens, all bonded at 109.5° from the other, and all with precisely the same bond angle. Each carbon-hydrogen bond is slightly polar (hydrogen has an electronegativity of 2.1, carbon 2.5), but because of this symmetry, the polarities cancel each other out and overall, methane is a non-polar molecule. The distinction is between Bond Polarity and Molecular polarity. The total polarity of a molecule is measured as Dipole Moment. The actual calculation of dipole moment isn't really necessary so much as an understanding of what it means. Frequently, a guesstimate of dipole moment is pretty easy once you understand the concept and until you get into the more advanced organic chemistry, exact values are of little value. Basically, the molecular polarity is, essentially, the summation of the vectors of all of the bond polarities in a molecule. Van der Waals Bonding. Van der Waals bonding is the collective name for three types of interactions: A Dipole is caused by an atom or molecule fragment having a higher electronegativity (this is a measure of its effective nuclear charge, and thus the attraction of the nucleus by electrons) than one to which it is attached. This means that it pulls electrons closer to it, and has a higher share of the electrons in the bond. Dipoles can cancel out by symmetry, eg: Carbon dioxide (O=C=O) is linear so there is no dipole, but the charge distribution is asymmetric causing a quadrupole moment (this acts similarly to a dipole, but is much weaker). Organometallic Compounds and Bonding. Organometallic chemistry combines aspects of inorganic chemistry and organic chemistry, because organometallic compounds are chemical compounds containing bonds between carbon and a metal or metalloid element. Organometallic bonds are different from other bonds in that they are not either truly covalent or truly ionic, but each type of metal has individual bond character. Cuprate (copper) compounds, for example, behave quite differently than Grignard reagents (magnesium), and so beginning organic chemists should concentrate on how to use the most basic compounds mechanistically, while leaving the explanation of exactly what occurs at the molecular level until later and more in-depth studies in the subject. Basic organometallic interactions are discussed fully in a later chapter. 

Resonance. Resonance refers to structures that are not easily represented by a single electron dot structure but that are intermediates between two or more drawn structures. Resonance is easily misunderstood in part because of the way certain chemistry textbooks attempt to explain the concept. In science, analogies can provide an aid to understanding, but analogies should not be taken too literally. It is sometimes best to use analogies to introduce a topic, but then explain the differences and inevitable complications as further details on a complicated subject. This is the case for resonance. Just as entropic principles cannot be applied to individual molecules, it is impossible to say whether or not any given individual molecule with a resonance structure is literally in one configuration or another. The actual situation on the molecular scale is that each configuration of the molecule contributes a percentage to the possible configurations, resulting in a "blend" of the possible structures. Changes in molecular shape occur so rapidly, and on such a tiny scale, that the actual physical locations of individual electrons cannot be precisely known (due to Heisenberg's Uncertainty Principle). The result of all that complexity is simply this: molecules with resonance structures are treated as mixtures of their multiple forms, with a greater percentage of probability given to the most stable configurations. The nuclei of the atoms are not moving when they are represented by resonance structure drawings. Rather, the electrons are portrayed as if they were moving instead. The true situation is that no one can say for certain exactly where any individual electron is at any specific moment, but rather electron location can be expressed as a probability only. What a dot structure is actually showing is where electrons almost certainly are located, therefore resonance structures indicate a split in those same probabilities. Chemists are absolutely certain where electrons are located when one carbon bonds four hydrogens (methane), but it is less certain where precisely any given electron is located when six carbons bond six hydrogens in a ring structrue (benzene). Resonance is an expression of this uncertainty, and is therefore the average of probable locations. Resonance structures are stabilizing in molecules because they allow electrons to lengthen their wavelengths and thereby lower their energy. This is the reason that benzene (C6H6) has a lower heat of formation than organic chemists would predict, not accounting for resonance. Other aromatic molecules have a similar stability, which leads to an overall entropic preference for aromaticity (a subject that will be covered fully in a later chapter). Resonance stability plays a major role in organic chemistry due to resonant molecules' lower energy of formation, so students of organic chemistry should understand this effect and practice spotting molecules stabilized by resonant forms. &lt;br&gt; In the Lewis structures above, carbonate (CO32-) has a resonance structure. Using laboratory procedures to measure the bond length of each bond, we do not find that one bond is shorter than the two others (remember, double bonds are shorter than single bonds), but instead that all bonds are of the same length somewhere between the length of typical double and single bonds. Resonance Structures. Resonance structures are diagrammatic tools used predominately in organic chemistry to symbolize resonant bonds between atoms in molecules. The electron density of these bonds is spread over the molecule, also known as the delocalization of electrons. Resonance contributors for the same molecule all have the same chemical formula and same sigma framework, but the pi electrons will be distributed differently among the atoms. Because Lewis dot diagrams often cannot represent the true electronic structure of a molecule, resonance structures are often employed to approximate the true electronic structure. Resonance structures of the same molecule are connected with a double-headed arrow. While organic chemists use resonance structures frequently, they are also used in inorganic structures, with nitrate as an example. Key characteristics. The key elements of resonance are: What resonance is not. Significantly, resonance structures do not represent different, isolatable structures or compounds. In the case of benzene, for example, there are two important resonance structures - which can be thought of as cyclohexa-1,3,5-trienes. There are other resonance forms possible, but because they are higher in energy than the triene structures (due to charge separation or other effects) they are less important and contribute less to the "real" electronic structure (average hybrid). However, this does not mean there are two different, interconvertable forms of benzene; rather, the true electronic structure of benzene is an average of the two structures. The six carbon-carbon bond lengths are identical when measured, which would be invalid for the cyclic triene. Resonance should also not be confused with a chemical equilibrium or tautomerism which are equilibria between compounds that have different sigma bonding patterns. Hyperconjugation is a special case of resonance. History. The concept of resonance was introduced by Linus Pauling in 1928. He was inspired by the quantum mechanical treatment of the H2+ ion in which an electron is located between two hydrogen nuclei. The alternative term mesomerism popular in German and French publications with the same meaning was introduced by Christopher Ingold in 1938 but did not catch on in the English literature. The current concept of Mesomeric effect has taken on a related but different meaning. The double headed arrow was introduced by the German chemist Arndt (also responsible for the Arndt-Eistert synthesis) who preferred the German phrase "zwischenstufe" or "intermediate phase". Due to confusion with the physical meaning of the word "resonance", as no elements do actually appear to be resonating, it is suggested to abandon the term "resonance" in favor of "delocalization" . Resonance energy would become delocalization energy and a "resonance structure" becomes contributing structure. The double headed arrows would get replaced by commas. Examples. The ozone molecule is represented by two resonance structures in the top of "scheme 2". In reality the two terminal oxygen atoms are equivalent and the hybrid structure is drawn on the right with a charge of -1/2 on both oxygen atoms and partial double bonds. The concept of benzene as a hybrid of two conventional structures (middle "scheme 2") was a major breakthrough in chemistry made by Kekule, and the two forms of the ring which together represent the total resonance of the system are called "Kekule structures". In the hybrid structure on the right the circle replaces three double bonds. The allyl cation (bottom "scheme 2") has two resonance forms and in the hybrid structure the positive charge is delocalized over the terminal methylene groups. 

Arrhenius Definition: Hydroxide and Hydronium Ions. The first and earliest definition of acids and bases was proposed in the 1800s by Swedish scientist Svante Arrhenius, who said that an acid was anything that dissolved in water to yield H+ ions (like stomach acid HCl, hydrochloric acid), and a base was anything that dissolved in water to give up OH- ions (like soda lye NaOH, sodium hydroxide). Acids and bases were already widely used in various occupations and activities of the time, so Arrhenius' definition merely attempted to explained well-known and long-observed phenomenon. Although simple, at the time this definition of the two types of substances was significant. It allowed chemists to explain certain reactions as ion chemistry, and it also expanded the ability of scientists of the time to predict certain chemical reactions. The definition left a great deal wanting, however, in that many types of reactions that did not involve hydroxide or hydronium ions directly remained unexplained. Many general chemistry classes (especially in the lower grades or introductory levels) still use this simple definition of acids and bases today, but modern organic chemists make further distinctions between acids and bases than the distinctions provided under Arrhenius's definition. Brønsted-Lowry Acids and Bases: Proton donors and acceptors. A new definition for acids and bases, building upon the one already proposed by Arrhenius, was brought forth independently by Johannes Nicolaus Brønsted and Thomas Martin Lowry in 1923. The new definition did not depend on a substance's dissolution in water for definition, but instead suggested that a substance was acidic if it readily donated a proton (H+) to a reaction and a substance was basic if it accepted a proton in a reaction. The major advantage of the updated definition was that it was not limited to aqueous solution. This definition of acids and bases allowed chemists to explain a great number of reactions that took place in protic or aprotic solvents that were not water, and it also allowed for gaseous and solid phase reactions (although those reactions are more rare). For example, the hypothetical acid HA will disassociate into H+ and A-: The Brønsted-Lowry definition of acids and bases is one of two definitions still in common use by modern chemists. Lewis Acids and Bases: Electron donors and acceptors. The second definition in widespread use deals not with a molecule's propensity for accepting or donating protons but rather with accepting or donating electrons, thereby demonstrating a slightly different emphasis and further broadening the explanatory and predictive powers of acid-base chemistry. Probably the most important aspect of Lewis acids and bases is which types of atoms can donate electrons, and which types of atoms can receive them. Essentially atoms with lone pairs, i.e. unshared pairs of electrons in an outer shell, have the capability of using those lone pairs to attract electron-deficient atoms or ions. This is why ammonia can bond a fourth hydrogen ion to create the ammonium ion; its lone pair of electrons can attract and bond to a free H+ ion in solution and hold on to it. For the same reason, methane cannot become methanium ion under ordinary circumstances, because the carbon in methane does not have any unshared pairs of electrons orbiting its nucleus. Generally speaking, Lewis acid are in the nitrogen, oxygen or halogen groups of the periodic table. Nucleophiles and Electrophiles. Whether or not an atom can donate or accept electrons it can be called a nucleophile or electrophile, respectively. Electrophiles (literally, "lovers of electrons") are attracted to electrons. Electrophiles therefore seek to pair with unshared electrons of other atoms. Nucleophiles, or "nucleus lovers", seek positively charged nuclei such as those available in acidic solutions as hydronium ions. It is important to note that electrophiles and nucleophiles are often ions, but sometimes they are not. Understanding electrophiles and nucleophiles goes beyond simply ideas of acids and bases. They are, in a majority of cases, the major players in organic reactions. As we will, over and over again, find reactions that are the result of nucleophiles "attacking" electrophiles. Keep in mind that the idea of nucleophiles and electrophiles is very related to the ideas of acids and bases in the Lewis context. But it is also important to understand that, while they are related, they are not exactly the same thing either. An ion or molecule can be a strong nucleophile and a weak base (e.g. N3-, RS-, I-, Br- and CN-). Another ion can be a poor nucleophile and a strong base ((CH3)3CO-, R2N-). And yet others are strong nucleophiles and strong bases (R3C-, RO-, HO-) and poor nucleophiles and poor bases (RCO2-, ROH, NH3). This will all be discussed in greater detail as the topics of specific reactions and reaction mechanisms are covered. In the meantime, try to bear in mind that nucleophiles are basic and electrophiles are acidic. pKa and Acidity. The acid dissociation constant of a substance is commonly called its pKa, and it is a measure of the negative log of the K value of an acid dissociation reaction. (The K value refers to the equilibrium calculations you learned how to perform in general chemistry—if you have forgotten your K's and Q's, now would be a good time to refresh your memory on the topic.) The lower the pKa value is, the more acidic (and consequently, less basic) a substance is. There is also a pKb value for all relevant substances, but it is common in organic chemistry to use pKa exclusively, even when discussing bases. This is because extremely high pKa values correlate exactly to extremely low pKb values, so there is no need to use both kinds of measurements. Any pKa value higher than seven means that a substance is not acidic when placed in water, but it does not mean that substance cannot be an acid. Alcohols are a good example of this: they can donate a hydrogen ion in chemical reactions but they do not do so readily, which makes them acidic but only very weakly so. Many of the acids in organic chemistry are considerably weaker than acids used for inorganic chemistry, so discussion of acid-base chemistry in organic reactions may not necessarily relate well to your previous understanding of the topic. 

Electron Dot Structures. Electron dot structures, also called "Lewis structures", give a representation of the valence electrons surrounding an atom. Each valence electron is represented by one dot, thus, a lone atom of hydrogen would be drawn as an "H" with one dot, whereas a lone atom of Helium would be drawn as an "He" with two dots, and so forth. Representing two atoms joined by a covalent bond is done by drawing the atomic symbols near to each other, and drawing a single line to represent a shared pair of electrons. It is important to note: a single valence electron is represented by a dot, whereas a pair of electrons is represented by a line. The covalent compound hydrogen fluoride, for example, would be represented by the symbol "H" joined to the symbol "F" by a single line, with three pairs (six more dots) surrounding the symbol "F". The line represents the two electrons shared by both hydrogen and fluorine, whereas the six paired dots represent fluorine's remaining six valence electrons. Dot structures are useful in illustrating simple covalent molecules, but the limitations of dot structures become obvious when diagramming even relatively simple organic molecules. The dot structures have no ability to represent the actual physical orientation of molecules, and they become overly cumbersome when more than three or four atoms are represented. Lewis dot structures are useful for introducing the idea of covalence and bonding in small molecules, but other model types have much more capability to communicate chemistry concepts. Drawing electron dot structures. &lt;br&gt; Some examples of electron dot structures for a few commonly encountered molecules from inorganic chemistry. A note about Gilbert N. Lewis. Lewis was born in Weymouth, Massachusetts as the son of a Dartmouth-graduated lawyer/broker. He attended the University of Nebraska at age 14, then three years later transferred to Harvard. After showing an initial interest in Economics, Gilbert Newton Lewis earned first a B.A. in Chemistry, and then a Ph.D. in Chemistry in 1899. For a few years after obtaining his doctorate, Lewis worked and studied both in the United States and abroad (including Germany and the Philippines) and he was even a professor at M.I.T. from 1907 until 1911. He then went on to U.C. Berkeley in order to be Dean of the College of Chemistry in 1912. In 1916 Dr. Lewis formulated the idea that a covalent bond consisted of a shared pair of electrons. His ideas on chemical bonding were expanded upon by Irving Langmuir and became the inspiration for the studies on the nature of the chemical bond by Linus Pauling. In 1923, he formulated the electron-pair theory of acid-base reactions. In the so-called Lewis theory of acids and bases, a "Lewis acid" is an electron-pair acceptor and a "Lewis base" is an electron-pair donor. In 1926, he coined the term "photon" for the smallest unit of radiant energy. Lewis was also the first to produce a pure sample of deuterium oxide (heavy water) in 1933. By accelerating deuterons (deuterium nuclei) in Ernest O. Lawrence's cyclotron, he was able to study many of the properties of atomic nuclei. During his career he published on many other subjects, and he died at age 70 of a heart attack while working in his laboratory in Berkeley. He had one daughter and two sons; both of his sons became chemistry professors themselves. Formal Charge. The formal charge of an atom is the charge that it would have if every bond were 100% covalent (non-polar). Formal charges are computed by using a set of rules and are useful for accounting for the electrons when writing a reaction mechanism, but they don't have any intrinsic physical meaning. They may also be used for qualitative comparisons between different resonance structures (see below) of the same molecule, and often have the same sign as the partial charge of the atom, but there are exceptions. The formal charge of an atom is computed as the difference between the number of valence electrons that a neutral atom would have and the number of electrons that "belong" to it in the Lewis structure when one counts lone pair electrons as belonging fully to the atom, while electrons in covalent bonds are split equally between the atoms involved in the bond. The total of the formal charges on an ion should be equal to the charge on the ion, and the total of the formal charges on a neutral molecule should be equal to zero. For example, in the hydronium ion, H3O+, the oxygen atom has 5 electrons for the purpose of computing the formal charge—2 from one lone pair, and 3 from the covalent bonds with the hydrogen atoms. The other 3 electrons in the covalent bonds are counted as belonging to the hydrogen atoms (one each). A neutral oxygen atom has 6 valence electrons (due to its position in group 16 of the periodic table); therefore the formal charge on the oxygen atom is 6 – 5 = +1. A neutral hydrogen atom has one electron. Since each of the hydrogen atoms in the hydronium atom has one electron from a covalent bond, the formal charge on the hydrogen atoms is zero. The sum of the formal charges is +1, which matches the total charge of the ion. In chemistry, a formal charge (FC) on an atom in a molecule is defined as: When determining the correct Lewis structure (or predominant resonance structure) for a molecule, the structure is chosen such that the formal charge on each of the atoms is minimized. 

The real heart of organic chemistry is the reactions. Everything that you study is geared to prepare you for organic syntheses and other chemical transformations performed in the lab. This chapter gives you the basic tools to begin looking at these reactions. Some basic reaction types. One way to organize organic reactions places them into a few basic categories: Other categories include: Sometimes one reaction can fall into more than one category. These classifications are just a tool and are not rigid. Addition reaction. Something is added to something else to produce a third thing. "Note: the letters A, B and C here represent any atomic, ionic or molecular species which can undergo this type of reaction." Elimination reaction. Something comes off of a molecule, resulting in two products. Substitution reactions. This involves the exchange of one group for another. Common reaction types include Rearrangement reactions. A molecule shifts or otherwise rearranges to form a different molecule. This typically happens when one molecule changes into an isomer of itself. 

Homolytic vs heterolytic cleavage. Two bonded atoms can come apart from each other in one of two ways. Either In homolytic cleavage, each atom leaves with one-half of the shared electrons (one electron for a single bond, or two for double bonds). A—B → A* + B* A* and B* represent uncharged radicals. The "*" represents an unbonded, unpaired valence electron. In heterolytic cleavage, one atom leaves with all of the previously shared electrons and the other atom gets none of them. A—B → A− + B+ "Homo" (from the Greek for same) indicates that each atom leaves with the same number of electrons from the bond. "Hetero" (from the Greek for different) refers to the fact that one atom gets all of the bonding electrons, while the other gets none. Polar reactions. Polar reactions occur when two bonded atoms come apart, one taking more of the shared electrons than the other. They involve heterolytic cleavage. The result is two charged species—one cation and one anion. Radical reactions. Radical reactions don't deal with charged particles but with radicals. Radicals are uncharged atoms or molecules with an incomplete octet of valence electrons. When a molecule comes apart by homolytic cleavage the result is two radicals. Although uncharged, radicals are usually very reactive because the unfilled octet is unstable and the radical can lower its energy by forming a bond in a way that allows it to fill its valence shell while avoiding any electrostatic charge.. 

« Introduction to reactions | Introduction to functional groups | Overview of Functional Groups » | Alkenes » 

« Foundational concepts | Synthesis of urea » Long ago, people observed the differences between compounds that were derived from living things and those that were not. There seemed to be an impassable gap between the properties of the two groups. Someone proposed the vital force theory to explain the difference. The theory said that there was a something called a vital force that dwelled within the organic material that did not exist in the nonorganic materials. However, to echo the words of President Ronald Reagan, it was "only" a theory. Synthesis of urea » 

« Vital force theory | Organic vs inorganic chemistry » 

« Synthesis of urea | Atomic structure » 

« Atomic structure | Shells and orbitals » Atoms are made up of a nucleus and electrons that orbit the nucleus. An atom in its natural, uncharged state has the same number of electrons as protons. If it gains or loses electrons, the atom is then referred to as an ion. The nucleus. The nucleus is made up of protons, which each have a positive charge, and neutrons, which have no charge. Neutrons and protons have about the same mass, and together account for most of the mass of the atom. Each of these particles is made up of even smaller particles, though the existence of these particles do not come into play at the energies and time spans in which most chemical reactions occur. Electrons. The electrons are negatively charged and fly around the nucleus of an atom at something like light speed. We cannot determine exactly where electrons are at any point in time, rather, we can only guess at the probability of finding an electron at a point in space relative to a nucleus at any point in time. The image depicts the Bohr model of the atom, in which the electrons inhabit discrete "orbitals" around the nucleus much like planets orbit the sun. Current models of the atomic structure hold that electrons occupy fuzzy clouds around the nucleus of specific shapes, some spherical, some dumbbell shaped, some with even more complex shapes. Even though the simpler Bohr model of atomic structure has been superseded, we still refer to these electron clouds as "orbitals". The number of electrons and the nature of the orbitals they occupy greatly influence the reactivity of atoms in organic chemistry. « Atomic structure | Nucleus and electrons | Shells and orbitals » 

« Nucleus and electrons | Filling electron shells » Electron orbitals. Electrons orbit atoms in clouds of distinct shapes and sizes. The electron clouds are layered one inside the other into units called shells (think nested Russian dolls), with the electrons occupying the smallest, innermost shell having the lowest energy state and the electrons in the largest, outermost shell having the highest energy state. The higher the energy state, the more potential energy the electron has, just like a rock at the top of a hill has more potential energy than a rock at the bottom of a valley. These concepts will be important in understanding later concepts like optical activity of chiral compounds as well as many interesting things outside the realm of organic chemistry (like how lasers work). Wave nature of electrons. Electrons behave as particles but also as waves. (Work by Albert Einstein and others revealed that in fact, light and all matter behaves with this dual nature, and it is most clearly observed in the tiniest particles.) One of the results of this observation is that electrons are not just in simple orbit around the nucleus as we imagine the moon to circle the earth, but instead occupy space as if they were a wave on the surface of a sphere. If you jump a jumprope you could imagine that the wave in the rope is in its fundamental frequency. The high and low points fall right in the middle, and the places where the rope doesn't move much (the nodes) occur only at the two ends. If you shake the rope fast enough in a rythmic way, using more energy than you would just jumping rope, you might be able to make the rope vibrate with a wavelength shorter than the fundamental. You them might see that the rope has more than one place along its length where it vibrates from its highest spot to its lowest spot. Furthermore, you'll see that there is one or more places (or nodes) along its length where the rope seems to move very little, if at all. Or consider stringed musical instruments. The sound made by these instruments comes from the different ways, or modes the strings can vibrate. We can refer to these different patterns or modes of vibrations as linear harmonics. Going from there, we can recognize that a drum makes sound by vibrations that occur across the 2-dimensional surface of the drumhead. Extending this now into three dimensions, we think of the electron as vibrating across a 3-dimensional sphere, and the patterns or modes of vibration are referred to as spherical harmonics. The mathematical analysis of spherical harmonics were worked out by the French mathematician Legendre long before anyone started to think about the shapes of electron orbitals. The algebraic expressions he developed, known as Legendre polynomials, describe the three dimension shapes of electron orbitals in much the same way that the expression x2+y2 = z describes a circle (or, for that matter, a drumhead). Many organic chemists need never actually work with these equations, but it helps to understand where the pictures we use to think about the shapes of these orbitals come from. Electron shells. Each different shell is subdivided into one or more orbitals, which also have different energy levels, although the energy difference between orbitals is less than the energy difference between shells. Longer wavelengths have less energy; the s orbital has the longest wavelength allowed for an electron orbiting a nucleus and this orbital is observed to have the lowest energy. Each sub-shell in the main electron shell has a characteristic shape, and are named by a letter. The sub-shells are: s, p, d, and f. As one progresses up through the shells (represented by the principle quantum number n) more types of orbitals become possible. S orbital. The s orbital is the orbital lowest in energy and is spherical in shape. Electrons in this orbital are in their fundamental frequency. P orbital. The next lowest-energy orbital is the p orbital. Its shape is often described as like that of a dumbbell. There are three p-orbitals each oriented along one of the 3-dimensional coordinates x, y or z. These three different p orbitals can be referred to as the px, py, and pz. The s and p orbitals are important for understanding most of organic chemistry as these are the orbitals that are occupied by the type of atoms that are most common in organic compounds. D orbital. There are 5 types of d orbitals. Three of them are roughly X-shaped, as shown here, and might be viewed as being shaped like a crossed pair of dumbbells . They are referred to as dxy, dxz&lt;/sub u&gt;, and dyz. Like the p-orbitals, these three d orbitals have a node at the origin of the coordinate system where the three axes all come together. Unlike the p orbitals, however, these three d orbitals are not oriented along the x, y, or z axes, but instead are oriented in between them. The dxy orbital, for instance, lies in the xy plane, but the lobes of the orbital point out in between the x and y axes. F orbital and beyond. There are 7 kinds of F orbitals, but we will not discuss their shapes here. F orbitals are filled in the elements of the lanthanide and actinide series, although electrons in these orbitals rarely come into play in organometallic reactions involving these elements. 

« Shells and orbitals | Octet rule and exceptions » When an atom or ion receives electrons into its orbitals, the orbitals and shells fill up in a particular manner. There are three principles that govern this process: 1) the Aufbau (build-up) principle, 2) the Pauli exclusion principle, and 3) Hund's rule. Exclusion principle. No more than one electron can have all four quantum numbers the same. What this translates to in terms of our pictures of orbitals is that each orbital can only hold two electrons, one "spin up" and one "spin down". Build-up principle. You may consider an atom as being "built up" from a naked nucleus by gradually adding to it one electron after another, until all the electrons it will hold have been added. Much as one fills up a container with liquid from the bottom up, so also are the orbitals of an atom filled from the lowest energy orbitals to the highest energy orbitals. However, the three p orbitals of a given shell all occur at the same energy level. So, how are they filled up? Is one of them filled full with the two electrons it can hold first? Or do each of the three orbitals receive one electron apiece before the any single orbital is double occupied. As it turns out, the latter situation occurs. Hund's rule. This rule is applicable only for those elements that have d electrons, and so is less important in organic chemistry (though it is important in organometallic chemistry). It says that filled and half-filled shells tend to have additional stability. In some instances, then, for example, the 4s orbitals will be filled before the 3d orbitals. Building atoms with quantum mechanics (advanced topic). The equations (like the Legendre polynomials that describe spherical harmonics, and thus the shapes of orbitals) of quantum mechanics are distinguished by four types of numbers. The first of these quantum numbers is referred to as the principal quantum number, and is indicated by n. This merely represents which shell electrons occupy, and shows up in the periodic chart as the rows of the periodic chart. It has integral values, n=1,2,3 . . . . The highest energy electrons of the atoms in the first row all have n=1. Those in the atoms in the second row all have n=2. As one gets into n=3, Hund's rule mixes it up a little bit, but when one gets to the end of the third row, at least, the electrons with the highest energy have n=3. The next quantum number is indicated by the letter m and indicates how many different types of shells an atom can have. Those elements in the first row can have just one, the s orbital. The elements in the second row can have two, the s and the p orbitals. The elements in the third row can have three, the s, p, and d orbitals. And so on. It may be funny to think of s, p and d as "numbers", but these are used as an historical and geometrical convenience. The third kind of quantum number ml specifies, for those kinds of orbitals that can have different shapes, which of the possible shapes one is referring to. So, for example, a 2pz orbital indicates three quantum numbers, represented respectively by the 2, the p and the z. Finally, the fourth quantum number is the spin of the electron. It has only two possible values, +1/2 or -1/2. Pretty much only computational chemists have to treat quantum numbers as numbers per se is equations. But it helps to know that the wide variety of elements of the periodic table and the different shapes and other properties of electron orbitals have a unifying principle--the proliferation of different shapes is not completely arbitrary, but is instead bounded by very specific rules. Afbau Principle (it means 'building up'):- It states that the orbitals should be filled according to their increasing energies Thus the lowest energy orbital which is available is filled first. The increasing order of energies of the various orbitals is:-  1s,2s,2p,3s,3p,4s,3d,4p,5s,4d,5p,6s,4f,5d,6p,7s,5f... The order of increasing of energy of orbitals can be calc. from(n+l) rule or 'Bohr bury rule' According to this rule, the value of n+l is the energy of the orbital and such on orbital will be filled up first. e.g. 4s orbital having lower value of(n+l) has lower energy than 3d orbital and hence 4s orbital is filled up first. For 4s orbital, n+l=4+0=4 For 3d orbital, n+l=3+2=5,therefore 4s orbital will be filled first. 





This book discusses proof-based linear algebra. The book was designed specifically for students who had not previously been exposed to mathematics as mathematicians view it. That is, as a subject whose goal is to "rigorously" prove theorems starting from clear consistent definitions. This book attempts to build students up from a background where mathematics is simply a tool that provides useful calculations to the point where the students have a grasp of the clear and precise nature of mathematics. For a milder introduction to linear algebra that is not too proof-based, see Introductory Linear Algebra wikibook. A more detailed discussion of the prerequisites and goal of this book is given in the introduction. 

What is calculus? Calculus is the broad area of mathematics dealing with such topics as instantaneous rates of change, areas under curves, and sequences and series. Underlying all of these topics is the concept of a limit, which consists of analyzing the behavior of a function at points ever closer to a particular point, but without ever actually reaching that point. As a typical application of the methods of calculus, consider a moving car. It is possible to create a function describing the displacement of the car (where it is located in relation to a reference point) at any point in time as well as a function describing the velocity (speed and direction of movement) of the car at any point in time. If the car were traveling at a constant velocity, then algebra would be sufficient to determine the position of the car at any time; if the velocity is unknown but still constant, the position of the car could be used (along with the time) to find the velocity. However, the velocity of a car cannot jump from zero to 35 miles per hour at the beginning of a trip, stay constant throughout, and then jump back to zero at the end. As the accelerator is pressed down, the velocity rises gradually, and usually not at a constant "rate" (i.e., the driver may push on the gas pedal harder at the beginning, in order to speed up). Describing such motion and finding velocities and distances at particular times cannot be done using methods taught in pre-calculus, whereas it is not only possible but straightforward with calculus. Calculus has two basic applications: "differential calculus" and "integral calculus". The simplest introduction to differential calculus involves an explicit series of numbers. Given the series (42, 43, 3, 18, 34), the differential of this series would be (1, -40, 15, 16). The new series is derived from the difference of successive numbers which gives rise to its name "differential". Rarely, if ever, are differentials used on an explicit series of numbers as done here. Instead, they are derived from a continuous function in a manner which is described later. Integral calculus, like differential calculus, can also be introduced via series of numbers. Notice that in the previous example, the original series can almost be derived solely from its differential. Instead of taking the difference, however, integration involves taking the sum. Given the first number of the original series, 42 in this case, the rest of the original series can be derived by adding each successive number in its differential (42+1, 43-40, 3+15, 18+16). Note that knowledge of the first number in the original series is crucial in deriving the integral. As with differentials, integration is performed on continuous functions rather than explicit series of numbers, but the concept is still the same. Integral calculus allows us to calculate the area under a curve of almost any shape; in the car example, this enables you to find the displacement of the car based on the velocity curve. This is because the area under the curve is the total distance moved, as we will soon see. Why learn calculus? Calculus is essential for many areas of science and engineering. Both make heavy use of mathematical functions to describe and predict physical phenomena that are subject to continuous change, and this requires the use of calculus. Take our car example: if you want to design cars, you need to know how to calculate forces, velocities, accelerations, and positions. All require calculus. Calculus is also necessary to study the motion of gases and particles, the interaction of forces, and the transfer of energy. It is also useful in business whenever rates are involved. For example, equations involving interest or supply and demand curves are grounded in the language of calculus. Calculus also provides important tools in understanding functions and has led to the development of new areas of mathematics including real and complex analysis, topology, and non-euclidean geometry. Notwithstanding calculus' "functional" utility (pun intended), many non-scientists and non-engineers have chosen to study calculus just for the challenge of doing so. A smaller number of persons undertake such a challenge and then discover that calculus is beautiful in and of itself. What is involved in learning calculus? Learning calculus, like much of mathematics, involves two parts: What you should know before using this text. There are some basic skills that you need before you can use this text. Continuing with our example of a moving car: Scope. The first four chapters of this textbook cover the topics taught in a typical high school or first year college course. The first chapter, ../Precalculus/, reviews those aspects of functions most essential to the mastery of calculus. The second, ../Limits/, introduces the concept of the limit process. It also discusses some applications of limits and proposes using limits to examine slope and area of functions. The next two chapters, ../Differentiation/ and ../Integration/, apply limits to calculate derivatives and integrals. The Fundamental Theorem of Calculus is used, as are the essential formulas for computation of derivatives and integrals without resorting to the limit process. The third and fourth chapters include articles that apply the concepts previously learned to calculating volumes, and as other important formulas. The remainder of the central calculus chapters cover topics taught in higher-level calculus topics: parametric and polar equations, sequences and series, multivariable calculus, and differential equations. The final chapters cover the same material, using formal notation. They introduce the material at a much faster pace, and cover many more theorems than the other two sections. They assume knowledge of some set theory and set notation. 

Functions are everywhere, from a simple correlation between distance and time to complex heat waves. This chapter focuses on the fundamentals of functions: the definition, basic concepts, and other defining aspects. It is very concept-heavy, and expect a lot of reading and understanding. However, this is simply a review and an introduction on what is to come in future chapters. Introduction. Whenever one quantity uniquely determines the value of another quantity, we have a function. That is, the set formula_1 uniquely determines the set formula_2. You can think of a "function" as a kind of machine. You feed the machine raw materials, and the machine changes the raw materials into a finished product. Think about dropping a ball from a bridge. At each moment in time, the ball is a height above the ground. The height of the ball is a function of time. It was the job of physicists to come up with a formula for this function. This type of function is called "real-valued" since the "finished product" is a number (or, more specifically, a real number). Think about a wind storm. At different places, the wind can be blowing in different directions with different intensities. The direction and intensity of the wind can be thought of as a function of position. This is a function of two real variables (a location is described by two values - an formula_3 and a formula_4) which results in a vector (which is something that can be used to hold a direction and an intensity). These functions are studied in multivariable calculus (which is usually studied after a one year college level calculus course). This a vector-valued function of two real variables. We will be looking at real-valued functions until studying multivariable calculus. Think of a real-valued function as an "input-output machine"; you give the function an input, and it gives you an output which is a number (more specifically, a real number). For example, the squaring function takes the input 4 and gives the output value 16. The same squaring function takes the input -1 and gives the output value 1. Notation. Functions are used so much that there is a special notation for them. The notation is somewhat ambiguous, so familiarity with it is important in order to understand the intention of an equation or formula. Though there are no strict rules for naming a function, it is standard practice to use the letters formula_5 , formula_6 , and formula_7 to denote functions, and the variable formula_3 to denote an independent variable. formula_4 is used for both dependent and independent variables. When discussing or working with a function formula_5 , it's important to know not only the function, but also its independent variable formula_3 . Thus, when referring to a function formula_5, you usually do not write formula_5, but instead formula_14 . The function is now referred to as "formula_5 of formula_3". The name of the function is adjacent to the independent variable (in parentheses). This is useful for indicating the value of the function at a particular value of the independent variable. For instance, if and if we want to use the value of formula_5 for formula_3 equal to formula_20 , then we would substitute 2 for formula_3 on both sides of the definition above and write This notation is more informative than leaving off the independent variable and writing simply 'formula_5' , but can be ambiguous since the parentheses next to formula_5 can be misinterpreted as multiplication, formula_25. To make sure nobody is too confused, follow this procedure: Description. There are many ways which people describe functions. In the examples above, a verbal description is given (the height of the ball above the earth as a function of time). Here is a list of ways to describe functions. The top three listed approaches to describing functions are the most popular. When a function is given a name (like in number 1 above), the name of the function is usually a single letter of the alphabet (such as formula_5 or formula_6). Some functions whose names are multiple letters (like the sine function formula_34 If we write formula_30 , then we know that How would we know the value of the function formula_5 at 3? We would have the following three thoughts: and we would write formula_46. The value of formula_5 at 3 is 11. Note that formula_48 means the value of the dependent variable when formula_3 takes on the value of 3. So we see that the number "11" is the output of the function when we give the number "3" as the input. People often summarize the work above by writing "the value of formula_5 at three is eleven", or simply "formula_5 of three equals eleven". Basic concepts of functions. The formal definition of a function states that a function is actually a "mapping" that associates the elements of one set called the domain of the function, formula_52, with the elements of another set called the "range" of the function, formula_53. For each value we select from the domain of the function, there exists "exactly one" corresponding element in the range of the function. The definition of the function tells us which element in the range corresponds to the element we picked from the domain. An example of how is given below. In mathematics, it is important to notice any repetition. If something seems to repeat, keep a note of that in the back of your mind somewhere. Here, notice that formula_54 and formula_55. Because formula_14 is equal to two different things, it must be the case that the other side of the equal sign to formula_14 is equal to the other. This property is known as the transitive property and could thus make the following equation below true:&lt;br&gt; Next, simplify — make your life easier rather than harder. In this instance, since formula_59 has formula_3 as a like-term, and the two terms are fractions added to the other, put it over a common denominator and simplify. Similar, since formula_61 is a mixed fraction, formula_62.&lt;br&gt; Multiply both sides by the reciprocal of the coefficient of formula_3, formula_68 from both sides by multiplying by it.&lt;br&gt; The value of formula_3 that makes formula_55 is formula_71.formula_75. Classically, the element picked from the domain is pictured as something that is fed into the function and the corresponding element in the range is pictured as the output. Since we "pick" the element in the domain whose corresponding element in the range we want to find, we have control over what element we pick and hence this element is also known as the "independent variable". The element mapped in the range is beyond our control and is "mapped to" by the function. This element is hence also known as the "dependent variable", for it depends on which independent variable we pick. Since the elementary idea of functions is better understood from the classical viewpoint, we shall use it hereafter. However, it is still important to remember the correct definition of functions at all times. Basic types of transformation. To make it simple, for the function formula_14, all of the possible formula_3 values constitute the domain, and all of the values formula_14 (formula_4 on the x-y plane) constitute the range. To put it in more formal terms, a function formula_5 is a mapping of some element formula_81, called the domain, to exactly one element formula_82, called the range, such that formula_83. The image below should help explain the modern definition of a function: The modern definition describes a function sufficiently such that it helps us determine whether some new type of "function" is indeed one. Now that each case is defined above, you can now prove whether functions are one of these function cases. Here is an example problem:  Notice that the only changing element in the function formula_5 is formula_3. To prove a function is one-to-one by applying the definition may be impossible because although two random elements of domain set formula_52 may not be many-to-one, there may be some elements formula_99 that may make the function many-to-one. To account for this possibility, we must prove that it is impossible to have some result like that. Assume there exists two distinct elements formula_115 that will result in formula_116. This would make the function many-to-one. In consequence,&lt;br&gt; formula_117 Subtract formula_89 from both sides of the equation. Subtract formula_120 from both sides of the equation. Factor formula_91 from both terms to the left of the equation. Multiply formula_124 to both sides of the equation. Add formula_126 to both sides of the equation. Notice that formula_128. However, this is impossible because formula_129 and formula_126 are distinct. Ergo, formula_131. No two distinct inputs can give the same output. Therefore, the function must be one-to-one. Basic concepts. There are a few more important ideas to learn from the modern definition of the function, and it comes from knowing the difference between domain, range, and codomain. We already discussed some of these, yet knowing about sets adds a new definition for each of the following ideas. Therefore, let us discuss these based on these new ideas. Let formula_52 and formula_53 be a set. If we were to combine these two sets, it would be defined as the "cartesian cross product" formula_134. The subset of this product is the function. The below definitions are a little confusing. The best way to interpret these is to draw an image. To the right of these definitions is the image that corresponds to it. Note that the codomain is not as important as the other two concepts. Take formula_135 for example:Because of the square root, the content in the square root should be strictly not smaller than 0.formula_136formula_137Thus the domain isformula_138 Correspondingly, the range will beformula_139 Other types of transformation. There is one more final aspect that needs to be defined. We already have a good idea of what makes a mapping a function (e.g. it cannot be one-to-many). However, three more definitions that you will often hear will be necessary to talk about: "injective", "surjective", "bijective". Again, the above definitions are often very confusing. Again, another image is provided to the right of the bullet points. Along with that, another example is also provided. Let us analyze the function formula_142.  Notice that the only changing element in the function formula_6 is formula_3. Let us check to see the conditions of this function. Is it injective? Let formula_145, and solve for formula_3. First, divide by formula_91. Then take the square root of formula_150. By definition, formula_151, so Notice, however, what we learned from the above manipulation. As a result of solving for formula_3, we found that there are two solutions for formula_3. However, this resulted in two different values from formula_156. This means that for some individual formula_3 that gives formula_4, there are two different inputs that result in the same value. Because formula_159 when formula_160, this is by definition non-injective. Is it surjective? As a natural consequence of what we learned about inputs, let us determine the range of the function. After all, the only way to determine if something is surjective is to see if the range applies to all real numbers. A good way to determine this is by finding a pattern involving our domains. Let us say we input a negative number for the domain: formula_161. This seemingly trivial exercise tells us that negative numbers give us non-negative numbers for our range. This is huge information! For formula_162, the function always results formula_163 for our range. For the set formula_53, we have elements in that set that have no mappings from the set formula_52. As such, formula_166 is the codomain of set formula_53. Therefore, this function is non-surjective! Tests for equations. The vertical line test. The vertical line test is a systematic test to find out if an equation involving formula_3 and formula_4 can serve as a function (with formula_3 the independent variable and formula_4 the dependent variable). Simply graph the equation and draw a vertical line through each point of the formula_3-axis. If any vertical line ever touches the graph at more than one point, then the equation is not a function; if the line always touches at most one point of the graph, then the equation is a function. The circle, on the right, is not a function because the vertical line intercepts two points on the graph when formula_173. The horizontal line and the algebraic 1-1 test. Similarly, the horizontal line test, though does not test if an equation is a function, tests if a function is injective (one-to-one). If any horizontal line ever touches the graph at more than one point, then the function is not one-to-one; if the line always touches at most one point on the graph, then the function is one-to-one. The algebraic 1-1 test is the non-geometric way to see if a function is one-to-one. The basic concept is that: Assume there is a function formula_5. If:formula_175, and formula_176, thenfunction formula_5 is one-to-one. Here is an example: prove that formula_178 is injective. Since the notation is the notation for a function, the equation is a function. So we only need to prove that it is injective. Let formula_91 and formula_89 be the inputs of the function and that formula_175. Thus, So, the result is formula_176, proving that the function is injective. Another example is proving that formula_190 is not injective. Using the same method, one can find that formula_191, which is not formula_176. So, the function is not injective. Remarks. The following arise as a direct consequence of the definition of functions: Functions are an important foundation of mathematics. This modern interpretation is a result of the hard work of the mathematicians of history. It was not until recently that the definition of the relation was introduced by René Descartes in "Geometry" (1637). Nearly a century later, the standard notation (formula_197) was first introduced by Leonhard Euler in "Introductio in Analysin Infinitorum" and "Institutiones Calculi Differentialis". The term function was also a new innovation during Euler's time as well, being introduced Gottfried Wilhelm Leibniz in a 1673 letter "to describe a quantity related to points of a curve, such as a coordinate or curve's slope." Finally, the modern definition of the function being the relation among sets was first introduced in 1908 by Godfrey Harold Hardy where there is a relation between two variables formula_3 and formula_4 such that "to some values of formula_3 at any rate correspond values of formula_4." For the person that wants to learn more about the history of the function, go to History of functions. Important functions. The functions listed below are essential to calculus. This section only serves as a review and scratches the surface of those functions. If there are any questions about those functions, please review the materials related to those functions before continuing. More about graphing will be explained in Chapter Polynomials. Polynomial functions are the most common and most convenient functions in the world of calculus. Their frequent appearances and their applications on computer calculations have made them important. Constant. When formula_202, the polynomial can be rewritten into the following function:formula_203, where formula_204 is a constant.The graph of this function is a horizontal line passing the point formula_205. Linear. When formula_206, the polynomial can be rewritten intoformula_207, where formula_208 are constants.The graph of this function is a straight line passing the point formula_209 and formula_210, and the slope of this function is formula_211. Quadratic. When formula_212, the polynomial can be rewritten intoformula_213, where formula_214 are constants.The graph of this function is a parabola, like the trajectory of a basketball thrown into the hoop. If there are questions about the quadratic formula and other basic polynomial concepts, please review the respective chapters in Algebra. Trigonometric. Trigonometric functions are also very important because it can connect algebra and geometry. Trigonometric functions are explained in detail here due to its importance and difficulty. Exponential and Logarithmic. Exponential and logarithmic functions have a unique identity when calculating the derivatives, so this is a great time to review those functions. A special number will be frequently seen in those functions: the Euler's constant, also known as the base of the natural logarithm. Notated as formula_215, it is defined as formula_216. Signum. The Signum (sign) function is simply defined asformula_217 Properties of functions. Sometimes, a lot of function manipulations cannot be achieved without some help from basic properties of functions. Domain and range. This concept is discussed above. Growth. Although it seems obvious to spot a function increasing or decreasing, without the help of graphing software, we need a mathematical way to spot whether the function is increasing or decreasing, or else our human minds cannot immediately comprehend the huge amount of information. Assume a function formula_14 with inputs formula_219, and that formula_220, formula_221, and formula_222 at all times.If for all formula_223 and formula_224, formula_225, then formula_14 is increasing in formula_227 If for all formula_223 and formula_224, formula_230, then formula_14 is decreasing in formula_227Example: In which intervals is formula_233 increasing? Firstly, the domain is important. Because the denominator cannot be 0, the domain for this function is formula_234 In formula_235, the growth of the function is:Let formula_236 and formula_222 Thus,formula_238formula_239 both formula_236 formula_241 formula_242 formula_239 formula_222 and formula_245 formula_241 formula_247So, formula_248formula_14 is decreasing in formula_235In formula_251Let formula_252 and formula_222 Thus,formula_238formula_239 both formula_236 formula_241formula_242 However, the sign of formula_259 in formula_251 cannot be determined. It can only be determined in formula_261.In formula_262 formula_239 formula_222 and formula_245 formula_241 formula_247 In formula_268 formula_269 formula_270As a result, formula_14 is decreasing in formula_262 and increasing in formula_268.In formula_274Let formula_275 and formula_222 Thus,formula_238formula_239 both formula_236 formula_241formula_242 formula_269 formula_270So, formula_284formula_14 is increasing in formula_274.Therefore, the intervals in which the function is increasing are formula_287. formula_75 After learning derivatives, there will be more ways to determine the growth of a function. Parity. The properties odd and even are associated with symmetry. While even functions have a symmetry about the formula_4-axis, odd functions are symmetric about the origin. In mathematical terms:A function is even when formula_290 A function is odd when formula_291Example: Prove that formula_233 is an even function. formula_293 formula_294 is an even function formula_295 Manipulating functions. Addition, Subtraction, Multiplication and Division of functions. For two real-valued functions, we can add the functions, multiply the functions, raised to a power, etc. If we add the functions formula_31 and formula_297 , we obtain formula_298 . If we subtract formula_31 from formula_297 , we obtain formula_301 . We can also write this as formula_302 . If we multiply the function formula_31 and the function formula_297 , we obtain formula_305 . We can also write this as formula_306 . If we divide the function formula_31 by the function formula_297 , we obtain formula_309 . If a math problem wants you to add two functions formula_5 and formula_6 , there are two ways that the problem will likely be worded: Similar statements can be made for subtraction, multiplication and division. Let formula_30 and: formula_190 . Let's add, subtract, multiply and divide. Composition of functions. We begin with a fun (and not too complicated) application of composition of functions before we talk about what composition of functions is. If we drop a ball from a bridge which is 20 meters above the ground, then the height of our ball above the earth is a function of time. The physicists tell us that if we measure time in seconds and distance in meters, then the formula for height in terms of time is formula_330 . Suppose we are tracking the ball with a camera and always want the ball to be in the center of our picture. Suppose we have formula_331 The angle will depend upon the height of the ball above the ground and the height above the ground depends upon time. So the angle will depend upon time. This can be written as formula_332 . We replace formula_7 with what it is equal to. This is the essence of composition. Composition of functions is another way to combine functions which is different from addition, subtraction, multiplication or division. The value of a function formula_5 depends upon the value of another variable formula_3 ; however, that variable could be equal to another function formula_6 , so its value depends on the value of a third variable. If this is the case, then the first variable is a function formula_7 of the third variable; this function (formula_7) is called the composition of the other two functions (formula_5 and formula_6). Let formula_30 and: formula_190 . The composition of formula_5 with formula_6 is read as either "f composed with g" or "f of g of x." Let Then Sometimes a math problem asks you compute formula_347 when they want you to compute formula_348 , Here, formula_7 is the composition of formula_5 and formula_6 and we write formula_352 . Note that composition is not commutative: Composition of functions is very common, mainly because functions themselves are common. For instance, squaring and sine are both functions: Thus, the expression formula_358 is a composition of functions: Since the function sine equals formula_360 if formula_361 , Since the function square equals formula_363 if formula_361 , Transformations. Transformations are a type of function manipulation that are very common. They consist of multiplying, dividing, adding or subtracting constants to either the input or the output. Multiplying by a constant is called dilation and adding a constant is called translation. Here are a few examples: Translations and dilations can be either horizontal or vertical. Examples of both vertical and horizontal translations can be seen at right. The red graphs represent functions in their 'original' state, the solid blue graphs have been translated (shifted) horizontally, and the dashed graphs have been translated vertically. Dilations are demonstrated in a similar fashion. The function has had its input doubled. One way to think about this is that now any change in the input will be doubled. If I add one to formula_3, I add two to the input of formula_5, so it will now change twice as quickly. Thus, this is a horizontal dilation by formula_373 because the distance to the formula_4-axis has been halved. A vertical dilation, such as is slightly more straightforward. In this case, you double the output of the function. The output represents the distance from the formula_3-axis, so in effect, you have made the graph of the function 'taller'. Here are a few basic examples where formula_91 is any positive constant: Inverse functions. We call formula_378 the inverse function of formula_14 if, for all formula_3 : A function formula_14 has an inverse function if and only if formula_14 is one-to-one. For example, the inverse of formula_384 is formula_385 . The function formula_135 has no inverse because it is not injective. Similarly, the inverse functions of trigonometric functions have to undergo transformations and limitations to be considered as valid functions. Notation. The inverse function of formula_5 is denoted as formula_388 . Thus, formula_388 is defined as the function that follows this rule To determine formula_388 when given a function formula_5 , substitute formula_388 for formula_3 and substitute formula_3 for formula_14 . Then solve for formula_388 , provided that it is also a function. Example: Given formula_398 , find formula_388 . Substitute formula_388 for formula_3 and substitute formula_3 for formula_14 . Then solve for formula_388 : To check your work, confirm that formula_409 :formula_410formula_411formula_412If formula_5 isn't one-to-one, then, as we said before, it doesn't have an inverse. Then this method will fail. Example: Given formula_414 , find formula_388. Substitute formula_388 for formula_3 and substitute formula_3 for formula_14 . Then solve for formula_388 : Since there are two possibilities for formula_388 , it's not a function. Thus formula_414 doesn't have an inverse. Of course, we could also have found this out from the graph by applying the Horizontal Line Test. It's useful, though, to have lots of ways to solve a problem, since in a specific case some of them might be very difficult while others might be easy. For example, we might only know an algebraic expression for formula_14 but not a graph. =External links= 

Math Tutorial -- Vectors. Figure 1: Displacement vectors in a plane. Vector formula_1 represents the displacement of George from Mary, while vector formula_2 represents the displacement of Paul from George. Vector formula_3 represents the displacement of Paul from Mary and formula_4. The quantities formula_5, formula_6, etc., represent the Cartesian components of the vectors. Before we can proceed further we need to explore the idea of a "vector". A vector is a quantity which expresses both magnitude and direction. Graphically we represent a vector as an arrow. In typeset notation a vector is represented by a boldface character, while in handwriting an arrow is drawn over the character representing the vector. Figure 1 shows some examples of "displacement vectors", i. e., vectors which represent the displacement of one object from another, and introduces the idea of vector addition. The tail of vector formula_2 is collocated with the head of vector formula_1, and the vector which stretches from the tail of formula_1 to the head of formula_2 is the sum of formula_1 and formula_2, called formula_3 in figure 1. Figure 2: Definition sketch for the angle formula_14 representing the orientation of a two dimensional vector. The quantities formula_5, formula_6, etc., represent the Cartesian components of the vectors in figure 2. A vector can be represented either by its Cartesian components, which are just the projections of the vector onto the Cartesian coordinate axes, or by its direction and magnitude. The direction of a vector in two dimensions is generally represented by the counterclockwise angle of the vector relative to the formula_17 axis, as shown in figure 2. Conversion from one form to the other is given by the equations where formula_20 is the magnitude of the vector. A vector magnitude is sometimes represented by absolute value notation: formula_21. Notice that the inverse tangent gives a result which is ambiguous relative to adding or subtracting integer multiples of formula_22. Thus the quadrant in which the angle lies must be resolved by independently examining the signs of formula_5 and formula_6 and choosing the appropriate value of formula_14. To add two vectors, formula_1 and formula_2, it is easiest to convert them to Cartesian component form. The components of the sum formula_4 are then just the sums of the components: Subtraction of vectors is done similarly, e. g., if formula_30, then A unit vector is a vector of unit length. One can always construct a unit vector from an ordinary vector by dividing the vector by its length: formula_32. This division operation is carried out by dividing each of the vector components by the number in the denominator. Alternatively, if the vector is expressed in terms of length and direction, the magnitude of the vector is divided by the denominator and the direction is unchanged. Unit vectors can be used to define a Cartesian coordinate system. Conventionally, formula_33, formula_34, and formula_35 indicate the formula_17, formula_37, and formula_38 axes of such a system. Note that formula_33, formula_34, and formula_35 are mutually perpendicular. Any vector can be represented in terms of unit vectors and its Cartesian components: formula_42. An alternate way to represent a vector is as list of components: formula_43. We tend to use the latter representation since it is somewhat more economical notation. There are two ways to multiply two vectors, yielding respectively what are known as the dot product and the cross product. The cross product yields another vector while the dot product yields a number. Here we will discuss only the dot product. Figure 3: Definition sketch for dot product. Given vectors formula_1 and formula_2, the dot product of the two is defined where formula_14 is the angle between the two vectors. An alternate expression for the dot product exists in terms of the Cartesian components of the vectors: It is easy to show that this is equivalent to the cosine form of the dot product when the formula_17 axis lies along one of the vectors, as in figure 3. Notice in particular that formula_50, while formula_51 and formula_52. Thus, formula_53 in this case, which is identical to the form given above. By the law of cosines we can also see that which is an alternate coordinate-free expression for the dot product. Figure 4: Definition figure for rotated coordinate system. The vector formula_55 has components formula_56 and formula_57 in the unprimed coordinate system and components formula_58 and formula_59 in the primed coordinate system. All that remains to be proven for equation (2.6) to hold in general is to show that it yields the same answer regardless of how the Cartesian coordinate system is oriented relative to the vectors. To do this, we must show that formula_60, where the primes indicate components in a coordinate system rotated from the original coordinate system. This can be shown nearly instantly by applying the pythagorean theorem. Due to the fact that R is invariant and represents the hypotenuse for both triangles (X, X', Y and Y') we can conclude: formula_61 Since the dot product can be written solely in terms of magnitudes, as we did above, if the magnitude of a vector is invariant the dot product of two vectors must also be invariant. To deduce a general formula for X' and Y' you will have to do a bit more thinking: Figure 2.4 shows the vector formula_55 resolved in two coordinate systems rotated with respect to each other. From this figure it is clear that formula_63. Focusing on the shaded triangles, we see that formula_64 and formula_65. Thus, we find formula_66. Similar reasoning shows that formula_67 (Just imagine to rotate the constructs in the image further 90° without changing the axis-names. You will instantly notice that in the second quadrant X is negative while Y positive). Thus, the new and old coordinates are related by This is true of the position vector. We can use it to extend the notion of vector to concepts other than position by stating that a pair of numbers is a vector "if and only if" its values change in exactly this way under rotation. Substituting this relation into our earlier expression for the dot product and using the trigonometric identity formula_69 results in which proves the complete equivalence of the two forms of the dot product quoted above. (Multiply out the above expression to verify this.) A numerical quantity which doesn't depend on which coordinate system is being used is called a scalar. The dot product of two vectors is a scalar. However, the components of a vector, taken individually, are not scalars, since the components change as the coordinate system changes. Since the laws of physics cannot depend on the choice of coordinate system being used, we insist that physical laws be expressed in terms of scalars and vectors, but not in terms of the components of vectors. In three dimensions the cosine form of the dot product remains the same, while the component form is 

Reflection and Refraction. Most of what we need to know about geometrical optics can be summarized in two rules, the laws of reflection and refraction. These rules may both be inferred by considering what happens when a plane wave segment impinges on a flat surface. If the surface is polished metal, the wave is "reflected", whereas if the surface is an interface between two transparent media with differing indices of refraction, the wave is partially reflected and partially "refracted". Reflection means that the wave is turned back into the half-space from which it came, while refraction means that it passes through the interface, acquiring a different direction of motion from that which it had before reaching the interface. &lt;br&gt; Figure 3.1 shows the wave vector and wave front of a wave being reflected from a plane mirror. The angles of incidence, formula_1, and reflection, formula_2, are defined to be the angles between the incoming and outgoing wave vectors respectively and the line normal to the mirror. The law of reflection states that formula_3. This is a consequence of the need for the incoming and outgoing wave fronts to be in phase with each other all along the mirror surface. This plus the equality of the incoming and outgoing wavelengths is sufficient to insure the above result. &lt;br&gt; Refraction, as illustrated in figure 3.2, is slightly more complicated. Since formula_4, the speed of light in the right-hand medium is less than in the left-hand medium. (Recall that the speed of light in a medium with refractive index formula_5 is formula_6.) The frequency of the wave packet doesn't change as it passes through the interface, so the wavelength of the light on the right side is less than the wavelength on the left side. Let us examine the triangle ABC in figure 3.2. The side AC is equal to the side BC times formula_7. However, AC is also equal to formula_8, or twice the wavelength of the wave to the left of the interface. Similar reasoning shows that formula_9, twice the wavelength to the right of the interface, equals BC times formula_10. Since the interval BC is common to both triangles, we easily see that Since formula_12 and formula_13 where formula_14 and formula_15 are the wave speeds to the left and right of the interface, formula_16 is the speed of light in a vacuum, and formula_17 is the (common) period, we can easily recast the above equation in the form This is called "Snell's law", and it governs how a ray of light bends as it passes through a discontinuity in the index of refraction. The angle formula_1 is called the incident angle and formula_2 is called the refracted angle. Notice that these angles are measured from the normal to the surface, not the tangent. Derivation for Law of Reflection. The derivation of Law of Reflection using Fermat's principle is straightforward. The Law of Reflection can be derived using elementary Calculus and Trigonometry. The generalization of the Law of Reflection is Snell's law, which is derived bellow using the same principle. The medium that light travels through doesn't change. In order to minimize the time for light travel between to points, we should minimize the path taken. 1. Total path length of the light is given by 2. Using Pythagorean theorem from Euclidean Geometry we see that 3. When we substitute both values of d1 and d2 for above, we get 4. In order to minimize the path traveled by light, we take the first derivative of L with respect to x. 5. Set both sides equal to each other. 6. We can now tell that the left side is nothing but formula_27 and the right side formula_28 means 7. Taking the inverse sine of both sides we see that the angle of incidence equals the angle of reflection Derivation for Snell's Law. The derivation of Snell's Law using Fermat's Priciple is straightforward. Snell's Law can be derived using elementary calculus and trigonometry. Snell's Law is the generalization of the above in that it does not require the medium to be the same everywhere. To mark the speed of light in different media refractive indices named n1 and n2 are used. Here formula_33 is the speed of light in the vacuum and formula_34 because all materials slow down light as it travels through them. 1. Time for the trip equals distance traveled divided by the speed. 2. Using the Pythagorean theorem from Euclidean Geometry we see that 3. Substituting this result into equation (1) we get 4. Differentiating and setting the derivative equal to zero gives 5. After careful examination the above equation we see that it is nothing but 6. Thus 7. Multiplying both sides by formula_42 we get 8. Substituting formula_44 for v1 and formula_45 for formula_46 we get 9. Simplifying both sides we get our final result 

The modern approach to relativity. Although the special theory of relativity was first proposed by Einstein in 1905, the modern approach to the theory depends upon the concept of a four-dimensional universe, that was first proposed by Hermann Minkowski in 1908. Minkowski's contribution appears complicated but is simply an extension of Pythagoras' Theorem: In 2 dimensions: formula_1 In 3 dimensions: formula_2 in 4 dimensions: formula_3 The modern approach uses the concept of invariance to explore the types of coordinate systems that are required to provide a full physical description of the location and extent of things. The modern theory of special relativity begins with the concept of "length". In everyday experience, it seems that the length of objects remains the same no matter how they are rotated or moved from place to place. We think that the simple length of a thing is "invariant". However, as is shown in the illustrations below, what we are actually suggesting is that length seems to be invariant in a three-dimensional coordinate system. The length of a thing in a two-dimensional coordinate system is given by Pythagoras's theorem: This two-dimensional length is not invariant if the thing is tilted out of the two-dimensional plane. In everyday life, a three-dimensional coordinate system seems to describe the length fully. The length is given by the three-dimensional version of Pythagoras's theorem: The derivation of this formula is shown in the illustration below. It seems that, provided all the directions in which a thing can be tilted or arranged are represented within a coordinate system, then the coordinate system can fully represent the length of a thing. However, it is clear that things may also be changed over a period of time. Time is another direction in which things can be arranged. This is shown in the following diagram: The length of a straight line between two events in space and time is called a "space-time interval". In 1908 Hermann Minkowski pointed out that if things could be rearranged in time, then the universe might be four-dimensional. He boldly suggested that Einstein's recently-discovered theory of Special Relativity was a consequence of this four-dimensional universe. He proposed that the space-time interval might be related to space and time by Pythagoras' theorem in four dimensions: Where "i" is the imaginary unit (sometimes imprecisely called formula_8), "c" is a constant, and "t" is the time interval spanned by the space-time interval, "s". The symbols "x", "y" and "z" represent displacements in space along the corresponding axes. In this equation, the 'second' becomes just another unit of length. In the same way as centimetres and inches are both units of length related by centimetres = 'conversion constant' times inches, metres and seconds are related by metres = 'conversion constant' times seconds. The conversion constant, "c" has a value of about 300,000,000 meters per second. Now formula_9 is equal to minus one, so the space-time interval is given by: Minkowski's use of the imaginary unit has been superseded by the use of advanced geometry that uses a tool known as the "metric tensor". The metric tensor permits the existence of "real" time and the negative sign in the expression for the square of the space-time interval originates in the way that distance changes with time when the curvature of spacetime is analysed (see advanced text). We now use real time but Minkowski's original equation for the square of the interval survives so that the space-time interval is still given by: Space-time intervals are difficult to imagine; they extend between one place and time and another place and time, so the velocity of the thing that travels along the interval is already determined for a given observer. If the universe is four-dimensional, then the space-time interval (rather than the spatial length) will be invariant. Whoever measures a particular space-time interval will get the same value, no matter how fast they are travelling. In physical terminology the invariance of the spacetime interval is a type of Lorentz Invariance. The invariance of the spacetime interval has some dramatic consequences. The first consequence is the prediction that if a thing is travelling at a velocity of "c" metres per second, then all observers, no matter how fast they are travelling, will measure the same velocity for the thing. The velocity "c" will be a universal constant. This is explained below. When an object is travelling at "c", the space time interval is zero, this is shown below: A space-time interval of zero only occurs when the velocity is "c" (if x&gt;0). All observers observe the same space-time interval so when observers observe something with a space-time interval of zero, they all observe it to have a velocity of "c", no matter how fast they are moving themselves. The universal constant, "c", is known for historical reasons as the "speed of light in a vacuum". In the first decade or two after the formulation of Minkowski's approach many physicists, although supporting Special Relativity, expected that light might not travel at exactly "c", but might travel at very nearly "c". There are now few physicists who believe that light in a vacuum does not propagate at "c". The second consequence of the invariance of the space-time interval is that clocks will appear to go slower on objects that are moving relative to you. Suppose there are two people, Bill and John, on separate planets that are moving away from each other. John draws a graph of Bill's motion through space and time. This is shown in the illustration below:  Being on planets, both Bill and John think they are stationary, and just moving through time. John spots that Bill is moving through what John calls space, as well as time, when Bill thinks he is moving through time alone. Bill would also draw the same conclusion about John's motion. To John, it is as if Bill's time axis is leaning over in the direction of travel and to Bill, it is as if John's time axis leans over. The space-time interval, formula_19, is invariant. It has the same value for all observers, no matter who measures it or how they are moving in a straight line. Bill's formula_20 equals John's formula_20 so: So, if John sees Bill measure a time interval of 1 second (formula_25) between two ticks of a clock that is at rest in Bill's frame, John will find that his own clock measures between these same ticks an interval formula_26, called coordinate time, which is greater than one second. It is said that clocks in motion slow down, relative to those on observers at rest. This is known as "relativistic time dilation of a moving clock". The time that is measured in the rest frame of the clock (in Bill's frame) is called the proper time of the clock. John will also observe measuring rods at rest on Bill's planet to be shorter than his own measuring rods, in the direction of motion. This is a prediction known as "relativistic length contraction of a moving rod". If the length of a rod at rest on Bill's planet is formula_27, then we call this quantity the proper length of the rod. The length formula_28 of that same rod as measured from John's planet, is called coordinate length, and given by This equation can be derived directly and validly from the time dilation result with the assumption that the speed of light is constant. The last consequence is that clocks will appear to be out of phase with each other along the length of a moving object. This means that if one observer sets up a line of clocks that are all synchronised so they all read the same time, then another observer who is moving along the line at high speed will see the clocks all reading different times. In other words observers who are moving relative to each other see different events as simultaneous. This effect is known as Relativistic Phase or the Relativity of Simultaneity. Relativistic phase is often overlooked by students of Special Relativity, but if it is understood then phenomena such as the twin paradox are easier to understand. The way that clocks go out of phase along the line of travel can be calculated from the concepts of the invariance of the space-time interval and length contraction. In the diagram above John is conventionally stationary. Distances between two points according to Bill are simple lengths in space (x) all at t=0 whereas John sees Bill's measurement of distance as a combination of a distance (X) and a time interval (T): Notice that the quantities represented by capital letters are proper lengths and times and in this example refer to John's measurements. Bill's distance, x, is the length that he would obtain for things that John believes to be X metres in length. For Bill it is John who has rods that contract in the direction of motion so Bill's determination "x" of John's distance "X" is given from: This relationship between proper and coordinate lengths was seen above to relate Bill's proper lengths to John's measurements. It also applies to how Bill observes John's proper lengths. Clocks that are synchronised for one observer go out of phase along the line of travel for another observer moving at formula_37 metres per second by :formula_38 seconds for every metre. This is one of the most important results of Special Relativity and should be thoroughly understood by students. The net effect of the four-dimensional universe is that observers who are in motion relative to you seem to have time coordinates that lean over in the direction of motion and consider things to be simultaneous that are not simultaneous for you. Spatial lengths in the direction of travel are shortened, because they tip upwards and downwards, relative to the time axis in the direction of travel, akin to a rotation out of three-dimensional space. Interpreting space-time diagrams. Great care is needed when interpreting space-time diagrams. Diagrams present data in two dimensions, and cannot show faithfully how, for instance, a zero length space-time interval appears. When diagrams are used to show both space and time it is important to be alert to space and time being related by Minkowski's equation and not by simple Euclidean geometry. The diagrams are only aids to understanding the approximate relation between space and time and it must not be assumed, for instance, that simple trigonometric relationships can be used to relate lines that represent spatial displacements and lines that represent temporal displacements. It is sometimes mistakenly held that the time dilation and length contraction results only apply for observers at x=0 and t=0. This is untrue. An inertial frame of reference is defined so that length and time comparisons can be made anywhere within a given reference frame. Time differences in one inertial reference frame can be compared with time differences anywhere in another inertial reference frame provided it is remembered that these differences apply to corresponding pairs of lines or pairs of planes of simultaneous events. Spacetime. In order to gain an understanding of both Galilean and Special Relativity it is important to begin thinking of space and time as being different dimensions of a four-dimensional vector space called spacetime. Actually, since we can't visualize four dimensions very well, it is easiest to start with only one space dimension and the time dimension. The figure shows a graph with time plotted on the vertical axis and the one space dimension plotted on the horizontal axis. An "event" is something that occurs at a particular time and a particular point in space. ("Julius X. wrecks his car in Lemitar, NM on 21 June at 6:17 PM.") A "world line" is a plot of the position of some object as a function of time (more properly, the time of the object as a function of position) on a spacetime diagram. Thus, a world line is really a line in spacetime, while an event is a point in spacetime. A horizontal line parallel to the position axis (x-axis) is a "line of simultaneity"; in Galilean Relativity all events on this line occur simultaneously for all observers. It will be seen that the line of simultaneity differs between Galilean and Special Relativity; in Special Relativity the line of simultaneity depends on the state of motion of the observer. In a spacetime diagram the slope of a world line has a special meaning. Notice that a vertical world line means that the object it represents does not move -- the velocity is zero. If the object moves to the right, then the world line tilts to the right, and the faster it moves, the more the world line tilts. Quantitatively, we say that Notice that this works for negative slopes and velocities as well as positive ones. If the object changes its velocity with time, then the world line is curved, and the instantaneous velocity at any time is the inverse of the slope of the tangent to the world line at that time. The hardest thing to realize about spacetime diagrams is that they represent the past, present, and future all in one diagram. Thus, spacetime diagrams don't change with time -- the evolution of physical systems is represented by looking at successive horizontal slices in the diagram at successive times. Spacetime diagrams represent the evolution of events, but they don't evolve themselves. The lightcone. Things that move at the speed of light in our four dimensional universe have surprising properties. If something travels at the speed of light along the x-axis and covers x meters from the origin in t seconds the space-time interval of its path is zero. formula_42 but formula_43 so: formula_44 Extending this result to the general case, if something travels at the speed of light in any direction into or out from the origin it has a space-time interval of 0: formula_45 This equation is known as the Minkowski Light Cone Equation. If light were travelling towards the origin then the Light Cone Equation would describe the position and time of emission of all those photons that could be at the origin at a particular instant. If light were travelling away from the origin the equation would describe the position of the photons emitted at a particular instant at any future time 't'. At the superficial level the light cone is easy to interpret. Its backward surface represents the path of light rays that strike a point observer at an instant and its forward surface represents the possible paths of rays emitted from the point observer. Things that travel along the surface of the light cone are said to be light- like and the path taken by such things is known as a null geodesic. Events that lie outside the cones are said to be space-like or, better still space separated because their space time interval from the observer has the same sign as space (positive according to the convention used here). Events that lie within the cones are said to be time-like or time separated because their space-time interval has the same sign as time. However, there is more to the light cone than the propagation of light. If the added assumption is made that the speed of light is the maximum possible velocity then events that are space separated cannot affect the observer directly. Events within the backward cone can have affected the observer so the backward cone is known as the "affective past" and the observer can affect events in the forward cone hence the forward cone is known as the "affective future". The assumption that the speed of light is the maximum velocity for all communications is neither inherent in nor required by four dimensional geometry although the speed of light is indeed the maximum velocity for objects if the principle of causality is to be preserved by physical theories (ie: that causes precede effects). The Lorentz transformation equations. The discussion so far has involved the comparison of interval measurements (time intervals and space intervals) between two observers. The observers might also want to compare more general sorts of measurement such as the time and position of a single event that is recorded by both of them. The equations that describe how each observer describes the other's recordings in this circumstance are known as the Lorentz Transformation Equations. (Note that the symbols below signify coordinates.) The table below shows the Lorentz Transformation Equations. See mathematical derivation of Lorentz transformation. Notice how the phase ( (v/c2)x ) is important and how these formulae for absolute time and position of a joint event differ from the formulae for intervals. A spacetime representation of the Lorentz Transformation.  Bill and John are moving at a relative velocity, v, and synchronise clocks when they pass each other. Both Bill and John observe an event along Bill's direction of motion. What times will Bill and John assign to the event? It was shown above that the relativistic phase was given by: formula_46. This means that Bill will observe an extra amount of time elapsing on John's time axis due to the position of the event. Taking phase into account and using the time dilation equation Bill is going to observe that the amount of time his own clocks measure can be compared with John's clocks using: formula_47. This relationship between the times of a common event between reference frames is known as the Lorentz Transformation Equation for time. Continue 

Transverse and Longitudinal Waves. With the exception of light, waves are undulations in some material medium. For instance, waves on a slinky are either "transverse", in that the motion of the material of the slinky is perpendicular to the orientation of the slinky, if you vibrate the slinky like a rope, or they are "longitudinal", with material motion in the direction of the stretched slinky, if you treat it like a spring. (See image on right) Ocean waves are simultaneously transverse and longitudinal, the net effect being (nearly) circular undulations in the position of water parcels. The oscillations in neighboring parcels are phased such that a "pattern" moves across the ocean surface. Some media support only longitudinal waves, others support only transverse waves, while yet others support both types. Sound waves are purely longitudinal in gases and liquids, but can be either type in solids. Mechanical transverse waves require a material medium and propogate by means of vibrations of the medium perpendicular to the direction of travel. Examples are water waves, ripples, seismic shear waves, and waves in stretched strings as above. Electromagnetic (EM) waves (such as light) are also transverse waves but they do not require a medium and thus can pass through a vacuum (see intro). They consist of oscillating electric (E) and magnetic (B) fields which are perpendicular to the direction of propagation while also being mutually perpendicular. EM waves are a disturbance of space itself, which can be thought of as being stretched and therefore being elastic and having a tension. The B and E fields are in phase as shown on the left. The fundamental S.I. unit of length (meter) is defined in terms of the speed of light in vacuum and the definition of the unit of time, the second. The previous definition was in terms of the wavelength of a particular color of light in the line spectrum of Krypton 86. The modern definition is more accurate as well as being the same to all observers regardless of their relative velocity. Longitudinal waves propogate by means of vibrations or disturbances in the medium that are in the same direction that the wave travels. Examples are sound (as above), seismic shock waves, slinky springs and part of the motion in ocean waves. Sound waves are a series of high-pressure compressions and low-pressure rarefaction (for this reason, sound is sometimes called a "pressure wave"). While it might be convenient to think of individual molecules vibrating back and forth around an equilibrium position (producing areas of high and low pressure), the individual molecules in a gas generally move randomly, and it is only in large numbers that this pattern is visible. In longitudinal waves, the convention is to describe displacement in the direction the wave is going (the direction of propagation) as positive, and displacement against that direction as negative. The pressure and displacement of a molecule at a point are formula_1 out of phase, so that, for example, when a molecule is farthest displaced it is in a normal pressure area, while molecules in compressions or rarefactions have displacement close to zero. Torsional waves consist of a twisting disturbance moving through a medium such as a wire or a rod. 

Sine Waves. A particularly simple kind of wave, the sine wave, is illustrated in figure 1.2. This has the mathematical form: where Figure 1.2: Definition sketch for a sine wave, showing the wavelength λ and the amplitude formula_2 and the phase φ at various points. So far we have only considered a sine wave as it appears at a particular time. All interesting waves move with time. The movement of a sine wave to the right, a distance formula_4 may be accounted for by replacing formula_5 in the above formula by formula_6. If this movement occurs in time formula_7, then the wave moves at velocity formula_8. Solving this for formula_4 and substituting yields a formula for the displacement of a sine wave as a function of both distance formula_5 and time formula_7: The time for a wave to move one wavelength is called the period of the wave: formula_13. Thus, we can also write Physicists actually like to write the equation for a sine wave in a slightly simpler form. Defining the wavenumber as formula_15 and the angular frequency as formula_16, we write We normally think of the frequency of oscillatory motion as the number of cycles completed per second and is given by formula_18. It is related to the angular frequency omega by formula_19. The angular frequency is used because it is directly analogous to the wavenumber, see above. Converting between the two is not difficult. Frequency is measured in units of hertz, abbreviated Hz; formula_20 and angular frequency formula_21 is in units of radians per second. The argument of the sine function is by definition an angle. We refer to this angle as the phase of the wave, formula_22. The difference in the phase of a wave at fixed time over a distance of one wavelength is formula_23, as is the difference in phase at fixed position over a time interval of one wave period. As previously noted, we call formula_2, the maximum displacement of the wave, the amplitude. Often we are interested in the intensity of a wave, which is defined as the square of the amplitude, formula_25. The wave speed we have defined above, formula_26, is actually called the phase speed. Since formula_27 and formula_28, we can write the phase speed in terms of the angular frequency and the wavenumber: 

Types of Waves. In order to make the above material more concrete, we now examine the characteristics of various types of waves which may be observed in the real world. Ocean Surface Waves.  &lt;br&gt; Figure 1.3: Wave on an ocean of depth formula_1. The wave is moving to the right and the particles of water at the surface move up and down as shown by the small vertical arrows. The particles move up, and gravity is the 'restoring' force. Water waves are also longitudinal, the water particles moving forward, then back. The restoring forces are more complex, but involve the inertia of the mass of water surrounding. Think of water waves as superimposed transverse and longitudinal waves. The net particle paths are nearly circular. This is typical of waves that travel along the boundary (interface) between two substances, in this case water and air. These waves are manifested as undulations of the ocean surface as seen in figure 1.3. The speed of ocean waves is given by the formula where formula_3 is a constant related to the strength of the Earth's gravity, formula_1 is the depth of the ocean, and the hyperbolic tangent is defined as 2.7  &lt;br&gt; Figure 1.4: Plot of the function formula_6. The dashed line shows our approximation formula_7 for formula_8. As figure 1.4 shows, for very small x, we can approximate the hyperbolic tangent by formula_7, while for very large x it is positive 1 for positive x and negative 1 for negative x. This leads to two limits: Since formula_10, the "shallow water" limit, which occurs when formula_11, yields a wave speed of while the "deep water" limit, which occurs when formula_13, yields Notice that the speed of shallow water waves depends only on the depth of the water and on formula_15. In other words, all shallow water waves move at the same speed. On the other hand, deep water waves of longer wavelength (and hence smaller wavenumber) move more rapidly than those with shorter wavelength. Waves for which the wave speed varies with wavelength are called "dispersive". Thus, deep water waves are dispersive, while shallow water waves are non-dispersive. For water waves with wavelengths of a few centimeters or less, surface tension becomes important to the dynamics of the waves. In the deep water case the wave speed at short wavelengths is actually given by the formula where the constant formula_17 is related to surface tension and depends on the surfaces involved. For an air-water interface near room temperature, formula_18. 

""' Sound Waves. Sound is a longitudinal compression-expansion wave through matter. The wave speed is where formula_2 and formula_3 are constants and formula_4 is the "absolute temperature". The absolute temperature is measured in Kelvins and is numerically given by where formula_6 is the temperature in Celsius degrees. The angular frequency of sound waves is thus given by The speed of sound in air at normal temperatures is about formula_8. ""' 

Light. Light moves in a vacuum at a speed of formula_1. In transparent materials it moves at a speed less than formula_2 by a factor formula_3 which is called the "refractive index" of the material: Often the refractive index takes the form where formula_6 is the wavenumber and formula_7 and formula_8 are constants characteristic of the material. The angular frequency of light in a transparent medium is thus so the frequency increases slightly with increasing "k". Typically, when "k" is near "k""R", the material becomes opaque. Ultimately, this is due to resonance between the light and the atoms of the materials. 

Superposition Principle. It is found empirically that as long as the amplitudes of waves in most media are small, two waves in the same physical location don't interact with each other. Thus, for example, two waves moving in the opposite direction simply pass through each other without their shapes or amplitudes being changed. When superimposed, the total wave displacement is just the sum of the displacements of the individual waves. This is called the "superposition principle". At sufficiently large amplitude the superposition principle often breaks down -- interacting waves may scatter off of each other, lose amplitude, or change their form. "Interference" is a consequence of the superposition principle. When two or more waves are superimposed, the net wave displacement is just the algebraic sum of the displacements of the individual waves. Since these displacements can be positive or negative, the net displacement can either be greater or less than the individual wave displacements. The former case is called "constructive interference", while the latter is called "destructive interference".  &lt;br&gt; Figure 1.5: Superposition (lower panel) of two sine waves (shown individually in the upper panel) with equal amplitudes and wavenumbers formula_1 and formula_2 Let us see what happens when we superimpose two sine waves with different wavenumbers. Figure 1.5 shows the superposition of two waves with wavenumbers formula_1 and formula_2. Notice that the result is a wave with about the same wavelength as the two initial waves, but which varies in amplitude depending on whether the two sine waves are in or out of phase. When the waves are in phase, constructive interference is occurring, while destructive interference occurs where the waves are out of phase.  &lt;br&gt; Figure 1.6: Superposition of two sine waves with equal amplitudes and wavenumbers formula_5 and formula_6 What happens when the wavenumbers of the two sine waves are changed? Figure 1.6 shows the result when formula_5 and formula_6. Notice that though the wavelength of the resultant wave is decreased, the locations where the amplitude is maximum have the same separation in formula_9 as in figure 1.5.  &lt;br&gt; Figure 1.7: Superposition of two sine waves with equal amplitudes and wavenumbers formula_5 and formula_11. If we superimpose waves with formula_5 and formula_11, as is shown in figure 1.7, we see that the formula_9 spacing of the regions of maximum amplitude has decreased by a factor of two. Thus, while the wavenumber of the resultant wave seems to be related to something like the "average" of the wavenumbers of the component waves, the spacing between regions of maximum wave amplitude appears to go inversely with the "difference" of the wavenumbers of the component waves. In other words, if formula_15 and formula_16 are close together, the amplitude maxima are far apart and vice versa.  &lt;br&gt; Figure 1.8: Representation of the wavenumbers and amplitudes of two superimposed sine waves. We can symbolically represent the sine waves that make up figures 1.5, 1.6, and 1.7 by a plot such as that shown in figure 1.8. The amplitudes and wavenumbers of each of the sine waves are indicated by vertical lines in this figure. The regions of large wave amplitude are called wave packets. Wave packets will play a central role in what is to follow, so it is important that we acquire a good understanding of them. The wave packets produced by only two sine waves are not well separated along the formula_9-axis. However, if we superimpose many waves, we can produce an isolated wave packet. For example, figure 1.9 shows the results of superimposing formula_18 sine waves with wavenumbers formula_19, formula_20, where the amplitudes of the waves are largest for wavenumbers near formula_21.  &lt;br&gt; Figure 1.9: Superposition of twenty sine waves with formula_22 and formula_23. In particular, we assume that the amplitude of each sine wave is proportional to formula_24, where formula_22 and formula_23. The amplitudes of each of the sine waves making up the wave packet in figure 1.9 are shown schematically in figure 1.10.  &lt;br&gt; Figure 1.10: Representation of the wavenumbers and amplitudes of 20 superimposed sine waves with formula_22 and formula_23. The quantity formula_29 controls the distribution of the sine waves being superimposed -- only those waves with a wavenumber formula_30 within approximately formula_29 of the central wavenumber formula_32 of the wave packet, i. e., for formula_33 in this case, contribute significantly to the sum. If formula_29 is changed to formula_35, so that wavenumbers in the range formula_36 contribute significantly, the wavepacket becomes narrower, as is shown in figures 1.11 and 1.12.  &lt;br&gt; Figure 1.11: Superposition of twenty sine waves with formula_22 and formula_38.  &lt;br&gt; Figure 1.12: Representation of the wavenumbers and amplitudes of 20 superimposed sine waves with formula_22 and formula_38. formula_29 is called the wavenumber spread of the wave packet, and it evidently plays a role similar to the difference in wavenumbers in the superposition of two sine waves -- the larger the wavenumber spread, the smaller the physical size of the wave packet. Furthermore, the wavenumber of the oscillations within the wave packet is given approximately by the central wavenumber. We can better understand how wave packets work by mathematically analyzing the simple case of the superposition of two sine waves. Let us define formula_42 where formula_15 and formula_16 are the wavenumbers of the component waves. Furthermore let us set formula_45. The quantities formula_32 and formula_29 are graphically illustrated in figure 1.8. We can write formula_48 and formula_49 and use the trigonometric identity formula_50 to find formula_51 formula_52 formula_53 (2.17) The sine factor on the bottom line of the above equation produces the oscillations within the wave packet, and as speculated earlier, this oscillation has a wavenumber formula_32 equal to the average of the wavenumbers of the component waves. The cosine factor modulates this wave with a spacing between regions of maximum amplitude of Thus, as we observed in the earlier examples, the length of the wave packet formula_56 is inversely related to the spread of the wavenumbers formula_29 (which in this case is just the difference between the two wavenumbers) of the component waves. This relationship is central to the uncertainty principle of quantum mechanics. 

Beats. Suppose two sound waves of slightly different frequencies impinge on your ear at the same time. The displacement perceived by your ear is the superposition of these two waves, with time dependence formula_1 (2.19) where we now have formula_2 and formula_3. What you actually hear is a tone with angular frequency formula_4 which fades in and out with period formula_5 (2.20) The "beat frequency" is simply formula_6 (2.21) Note how beats are the time analog of wave packets -- the mathematics are the same except that frequency replaces wavenumber and time replaces space. 

Interferometers. An interferometer is a device which splits a beam of light into two sub-beams, shifts the phase of one sub-beam with respect to the other, and then superimposes the sub-beams so that they interfere constructively or destructively, depending on the magnitude of the phase shift between them. In this section we study the Michelson interferometer and interferometric effects in thin films. The Michelson Interferometer.  &lt;br&gt; Figure 1.13: Sketch of a Michelson interferometer. The American physicist Albert Michelson invented the optical interferometer illustrated in figure 1.13. The incoming beam is split into two beams by the half-silvered mirror. Each sub-beam reflects off of another mirror which returns it to the half-silvered mirror, where the two sub-beams recombine as shown. One of the reflecting mirrors is movable by a sensitive micrometer device, allowing the path length of the corresponding sub-beam, and hence the phase relationship between the two sub-beams, to be altered. As figure 1.13 shows, the difference in path length between the two sub-beams is formula_1 because the horizontal sub-beam traverses the path twice. Thus, constructive interference occurs when this path difference is an integral number of wavelengths, i. e., where formula_3 is the wavelength of the light and formula_4 is an integer. Note that formula_4 is the number of wavelengths that fits evenly into the distance formula_1. 

Thin Films.  &lt;br&gt; Figure 1.14: Plane light wave normally incident on a transparent thin film of thickness formula_1 and index of refraction formula_2. Partial reflection occurs at the front surface of the film, resulting in beam A, and at the rear surface, resulting in beam B. Much of the wave passes completely through the film, as with C. One of the most revealing examples of interference occurs when light interacts with a thin film of transparent material such as a soap bubble. Figure 1.14 shows how a plane wave normally incident on the film is partially reflected by the front and rear surfaces. The waves reflected off the front and rear surfaces of the film interfere with each other. The interference can be either constructive or destructive depending on the phase difference between the two reflected waves. If the wavelength of the incoming wave is formula_3, one would naively expect constructive interference to occur between the A and B beams if formula_4 were an integral multiple of formula_3. Two factors complicate this picture. First, the wavelength inside the film is not formula_3, but formula_7, where formula_8 is the index of refraction of the film. Constructive interference would then occur if formula_9. Second, it turns out that an additional phase shift of half a wavelength occurs upon reflection when the wave is incident on material with a higher index of refraction than the medium in which the incident beam is immersed. This phase shift doesn't occur when light is reflected from a region with lower index of refraction than felt by the incident beam. Thus beam B doesn't acquire any additional phase shift upon reflection. As a consequence, constructive interference actually occurs when while destructive interference results when When we look at a soap bubble, we see bands of colors reflected back from a light source. What is the origin of these bands? Light from ordinary sources is generally a mixture of wavelengths ranging from roughly formula_12 (violet light) to formula_13 (red light). In between violet and red we also have blue, green, and yellow light, in that order. Because of the different wavelengths associated with different colors, it is clear that for a mixed light source we will have some colors interfering constructively while others interfere destructively. Those undergoing constructive interference will be visible in reflection, while those undergoing destructive interference will not. Another factor enters as well. If the light is not normally incident on the film, the difference in the distances traveled between beams reflected off of the front and rear faces of the film will not be just twice the thickness of the film. To understand this case quantitatively, we need the concept of refraction, which will be developed later in the context of geometrical optics. However, it should be clear that different wavelengths will undergo constructive interference for different angles of incidence of the incoming light. Different portions of the thin film will in general be viewed at different angles, and will therefore exhibit different colors under reflection, resulting in the colorful patterns normally seen in soap bubbles. 

Math Tutorial -- Derivatives.  &lt;br&gt; Figure 1.15: Estimation of the derivative, which is the slope of the tangent line. When point B approaches point A, the slope of the line AB approaches the slope of the tangent to the curve at point A. This section provides a quick introduction to the idea of the derivative. For a more detailed discussion and exploration of the differentiation and of Calculus, see Calculus and Differentiation. Often we are interested in the slope of a line tangent to a function formula_1 at some value of formula_2. This slope is called the derivative and is denoted formula_3. Since a tangent line to the function can be defined at any point formula_2, the derivative itself is a function of formula_2: As figure 1.15 illustrates, the slope of the tangent line at some point on the function may be approximated by the slope of a line connecting two points, A and B, set a finite distance apart on the curve: As B is moved closer to A, the approximation becomes better. In the limit when B moves infinitely close to A, it is exact. Table of Derivatives. Derivatives of some common functions are now given. In each case formula_8 is a constant. The product and chain rules are used to compute the derivatives of complex functions. For instance, and 

Group Velocity. We now ask the following question: How fast do wave packets move? Surprisingly, we often find that wave packets move at a speed very different from the phase speed, formula_1, of the wave composing the wave packet. We shall find that the speed of motion of wave packets, referred to as the group velocity, is given by Dispersive waves are waves in which the phase speed varies with wavenumber. It is easy to show that dispersive waves have unequal phase and group velocities, while these velocities are equal for non-dispersive waves. Derivation of Group Velocity Formula. We now derive equation (1.36). It is easiest to do this for the simplest wave packets, namely those constructed out of the superposition of just two sine waves. We will proceed by adding two waves with full space and time dependence: After algebraic and trigonometric manipulations familiar from earlier sections, we find where as before we have formula_11, formula_12, formula_13, and formula_14. Again think of this as a sine wave of frequency formula_15 and wavenumber formula_16 modulated by a cosine function. In this case the modulation pattern moves with a speed so as to keep the argument of the cosine function constant: Differentiating this with respect to formula_18 while holding formula_19 and formula_20 constant yields In the limit in which the deltas become very small, this reduces to the derivative which is the desired result. 

Examples. We now illustrate some examples of phase speed and group velocity by showing the displacement resulting from the superposition of two sine waves, as given by equation (1.38), in the formula_1-formula_2 plane. This is an example of a spacetime diagram, of which we will see many examples later on.  &lt;br&gt; Figure 1.16: Net displacement of the sum of two traveling sine waves plotted in the formula_3 plane. The short vertical lines indicate where the displacement is large and positive, while the short horizontal lines indicate where it is large and negative. One wave has formula_4 and formula_5, while the other has formula_6 and formula_7. Thus, formula_8 and formula_9 and we have formula_10. Notice that the phase speed for the first sine wave is formula_11 and for the second wave is formula_12. Thus, formula_13 in this case. Figure 1.16 shows a non-dispersive case in which the phase speed equals the group velocity. The regions with vertical and horizontal hatching (short vertical or horizontal lines) indicate where the wave displacement is large and positive or large and negative. Large displacements indicate the location of wave packets. The positions of waves and wave packets at any given time may therefore be determined by drawing a horizontal line across the graph at the desired time and examining the variations in wave displacement along this line. The crests of the waves are indicated by regions of short vertical lines. Notice that as time increases, the crests move to the right. This corresponds to the motion of the waves within the wave packets. Note also that the wave packets, i. e., the broad regions of large positive and negative amplitudes, move to the right with increasing time as well. Since velocity is distance moved formula_14 divided by elapsed time formula_15, the slope of a line in figure 1.16, formula_16, is one over the velocity of whatever that line represents. The slopes of lines representing crests (the slanted lines, not the short horizontal and vertical lines) are the same as the slopes of lines representing wave packets in this case, which indicates that the two move at the same velocity. Since the speed of movement of wave crests is the phase speed and the speed of movement of wave packets is the group velocity, the two velocities are equal and the non-dispersive nature of this case is confirmed.  &lt;br&gt; Figure 1.17: Net displacement of the sum of two traveling sine waves plotted in the formula_1-formula_2 plane. One wave has formula_19 and formula_5, while the other has formula_21 and formula_22. In this case formula_23 while formula_24, so the group velocity is formula_25. However, the phase speeds for the two waves are formula_26 and formula_27. The average of the two phase speeds is about formula_28, so the group velocity is about twice the average phase speed in this case.  &lt;br&gt; Figure 1.18: Net displacement of the sum of two traveling sine waves plotted in the formula_1-formula_2 plane. One wave has formula_4 and formula_7, while the other has formula_6 and formula_5. Can you figure out the group velocity and the average phase speed in this case? Do these velocities match the apparent phase and group speeds in the figure? Figure 1.17 shows a dispersive wave in which the group velocity is twice the phase speed, while figure 1.18 shows a case in which the group velocity is actually opposite in sign to the phase speed. See if you can confirm that the phase and group velocities seen in each figure correspond to the values for these quantities calculated from the specified frequencies and wavenumbers. 

Problems.  &lt;br&gt; Figure 1.19: Sketch of a police radar.  &lt;br&gt; Figure 1.20: Sketch of a Fabry-Perot interferometer.  &lt;br&gt; Figure 1.21: Sketch of a weird dispersion relation. 

The wave is a physical phenomenon that is found in a variety of contexts, a perturbation in the medium. You undoubtedly know about ocean waves and have probably played with a stretched slinky toy, producing undulations which move rapidly along the slinky. Other examples of waves are sound, vibrations in solids, and light (photons have the particularity of sharing characteristics of both waves and particles). Despite the vast differences in these different types of waves, they all consist of a fluctuation around an equilibrium position caused by a restoring force, and can be described by the same types of mathematical equations. Understanding these equations is a powerful tool because it will let you understand the basics of a wide variety of seemingly unrelated phenomena. The purpose of this section is to describe the "kinematics" of waves, i.e., to provide tools for describing the form and motion of all waves irrespective of their underlying physical mechanisms. In this chapter we learn first about the basic properties of waves and introduce the most common type of wave, called the sinusoidal (sine) wave. In fact just about any type of wave can be expressed as a combination of sine functions using a technique proposed by Joseph Fourier. Examples of waves seen in the real world are presented. We then learn about the superposition principle, which allows us to construct complex wave patterns by superimposing sine waves. Using these ideas, we discuss the related ideas of beats and interferometry. Finally, the ideas of wave packets and group velocity are introduced. 

=Methane= « Foundational concepts | Ethane » Methane, CH4 is the simplest organic molecule. It is a gas at standard temperature and pressure and is a carbon atom with four hydrogen atoms bonded to it in a tetrahedral shape. This is a flattened, two-dimensional representation of methane that you will see commonly. The true three-dimensional form of methane does not have any 90 degree angles between bonded hydrogens. The bonds point to the four corners of a tetrahedron, forming cos−1(−1/3) ≈ 109.5 degree bond angles. « Foundational concepts | « Alkanes | Methane | Ethane » | Introduction to reactions » 

=Ethane= « Methane | Propane thru decane » Ethane (/ˈɛθeɪn/ or /ˈiːθeɪn/) is a chemical compound, or hydrocarbon with a chemical formula of C2H6. At standard temperature and pressure, ethane is a colorless, odorless gas. Ethane is isolated on an industrial scale from natural gas, and as a byproduct of petroleum refining. Its chief use is as petrochemical feedstock for ethylene production. It could be described as two methane molecules attached to each other, less two hydrogens. Number of hydrogens to carbons. This equation describes the relationship between the number of hydrogen and carbon atoms in alkanes: where "C" and "H" are used to represent the number of carbon and hydrogen atoms present in one molecule. If C = 2, then H = 6. Many textbooks put this in the following format: where "CN" and "H2N+2" represent the number of carbon and hydrogen atoms present in one molecule. If CN = 3, then H2N+2 = 2(3) + 2 = 8. (For this formula look to the "N" for the number, the "C" and the "H" letters themselves do not change.) « Organic chemistry | « Alkanes | « Methane | Ethane | Propane thru decane » | Introduction to reactions » 

All help is welcome. Optics. "about this book" Contents. X-ray optics /Authors/ External Resources.  __NOEDITSECTION__ 

Wikify. Wikify is a program written by Josh Ferguson specifically designed to turn a standard text file into a html/wiki format. It is being used to do the color syntax highlighting for the C++ code in the C++ section of this textbook. Author's Statement: Wikify can not yet wikify itself, due to the fact that it has html codes in it, which confuses itself. I plan to rewrite it soon so it can wikify itself. Right now though, it works good with simple code. And it can be modified to work with other languages by changing the variables in it. Later on I'll probably move all the important variables into text files so it will be more easy to change things for other languages (i.e. by changing the text files, or extension checks), etc... The code was written really fast, so the code isn't perfect, even though the author stubbornly maintains that it is close. Code Snapshot as of 7-22-03 The escape sequences in the code give the wikiscript problems, so here is a link to it. Click on it, and it will take you to a page, where you can click on it, and it will come up. This should be compilable by any standard C++ compiler. 

Derivative BASIC. BASIC was developed in 1963 at Dartmouth College in Hanover, New Hampshire as a teaching language. The acronym BASIC stands for Beginner's All-Purpose Symbolic Instruction Code. In 1964, John G. Kemeny and Thomas E. Kurtz designed the original BASIC language at Dartmouth College in New Hampshire. Programming Preliminaries. Anybody can write a program. A background in mathematics or science is not required. Patience, practice, and an interest in the subject matter should suffice, along with the required software and hardware. Understanding programs can appear daunting at first, but their reliance on logical operations allow for easy learning of commands which you will commonly see in many programs. A program itself is merely a series of commands in the order in which they are to be executed. That is to say, that the first line is the beginning of the program! All programs a user uses from day to day, including browsers (Internet Explorer, Firefox, etc.) and operating systems (Windows, Linux and Mac OS, etc) are separate sets of lines of code, which aim to fulfill tasks. The amount of code is dependent on how simple the task generally, and different types of code may be used for the advantages they give. BASIC is considered an excellent starting point for moving onto other languages, and can be useful for simple programs. Programming Languages. Programming languages allow people to give instructions to a computer with commands that both the computer and the programmer can understand. Different programming languages use different commands and different rules for entering those commands; similar to the way people speak different words to each other with the same meaning. One person may say "hello", while another says "hola", which appear different but express the same thought. Computer programming languages can be similar to each other in the same way that human languages, such as French and Spanish, can be. Programming languages that are similar are usually referred to as related languages. Once a person learns a programming language, it is easier to then learn other programming languages, especially those related to the first one, as many similarities in structure are shared between languages, especially those with a common ancestor. The language taught here, BASIC, is easier to learn than others as its commands are similar to English and it has a simple set of rules for entering them. Program. A program is defined as an instruction set that describes the logical steps the computer will follow to solve a particular problem. With programming the user is able to understand and communicate with the computer. Basics of BASIC. "Section for chapters detailing the basics of BASIC; i.e. data types, control structures..." 

Pascal is an influential computer programming language named after the mathematician . It was invented by in 1968 as a research project into the nascent field of compiler theory. The backronym PASCAL standing for "primary algorithmic scientific commercial application language" highlights its suitability for computing tasks in science, making it certainly usable for general programming as well. Contents. Standard Pascal Extensions Appendix Alternative resources. Tutorials, Textbooks, and the like: References, Articles on certain topics: 

« Filling electron shells | Molecular orbitals » The octet rule refers to the tendency of atoms to prefer to have eight electrons in their valence shell. The main exception to the rule is hydrogen, which is at lowest energy when it has two electrons in its valence shell. Other notable exceptions are aluminum and boron, which can function well with six valence electrons; and some atoms beyond group three on the periodic table that can have over 8 electrons, including sulfur. Additionally, some noble gasses can form compounds when expanding their valence shell. The other tendency of atoms with regard to their electrons is to maintain a neutral charge. Only the noble gasses have zero charge with filled valence octets. All of the other elements have a charge when they have eight electrons all to themselves. The result of these two guiding principles is the explanation for much of the reactivity and bonding that is observed within atoms; atoms seeking to share electrons in a way that minimizes charge while fulfilling an octet in the valence shell. « Foundational concepts | « Filling electron shells | Molecular orbitals » 

« Acids and bases | Electron donors and acceptors » The first and earliest definition of acids and bases was proposed in the 1800s by Swedish scientist Svante Arrhenius, who said that an acid was anything that dissolved in water to give up H+ ions,Such as hydrochloric acid and a base was anything that dissolved in water to give up OH- ions such a sodium hydroxide. The Brønsted-Lowry definition of 1923 broadened this idea a bit: The focus of this definition is on donating and accepting protons, and is not limited to aqueous solution. The Brønsted-Lowry definition of acids and bases is one of two definitions we commonly use. The second definition deals not with protons but with electrons, and has a slightly different emphasis. « Foundational concepts | « Resonance | « Acids and bases | Proton donors and acceptors | Alkanes » 

In organic chemistry we look at the hybridization of electron orbitals in to something called molecular orbitals. The s and p orbitals in a carbon atom combine into four hybridized orbitals that repel each other in a shape much like that of four balloons tied together. Remember what a tetrahedron looks like ? I knew you did ... 

The goal for this book is to create a complete, free, open-content, well-organized online book for the D programming language. D is a programming language being designed as a successor to C++. Until this page gets better written and more informative, the D home can be found here. Introduction. This book aims at beginners learning D language. It will cover all the language basics and some design aspects. In addition it will introduce topics like multi-threading, GUI programming and standard library to get you started with real-world applications. Overview. To quote Walter Bright, the author of the D Programming Language: D is a general purpose systems and applications programming language. It is a higher level language than C++, but retains the ability to write high performance code and interface directly with the operating system API's and with hardware. D is well suited to writing medium to large scale million line programs with teams of developers. D is easy to learn, provides many capabilities to aid the programmer, and is well suited to aggressive compiler optimization technology. D is not a scripting language, nor an interpreted language. It doesn't come with a VM, a religion, or an overriding philosophy. It's a practical language for practical programmers who need to get the job done quickly, reliably, and leave behind maintainable, easy to understand code. D is the culmination of decades of experience implementing compilers for many diverse languages, and attempting to construct large projects using those languages. D draws inspiration from those other languages (most especially C++) and tempers it with experience and real world practicality. D is a statically-typed, multi-paradigm language supporting imperative programming, object-oriented programming, and template metaprogramming. It also supports generics and design by contract. Main features of D. D has many features not seen in C++, implementing design by contract, unit testing, true modules, automatic memory management (garbage collection), first class arrays, associative arrays, dynamic arrays, array slicing, nested functions, inner classes, closures (anonymous functions), and has a reengineered template syntax. D retains C++'s ability to do low-level coding, and adds to it with support for an integrated inline assembler. C++ multiple inheritance is replaced by single inheritance with interfaces and mixins. D's declaration, statement and expression syntax closely matches that of C++. The inline assembler typifies the differentiation between D and application languages like Java and C#. An inline assembler allows a programmer to enter machine-specific assembly code alongside standard D code—a technique often used by systems programmers to access the low-level features of the processor needed to run programs that interface directly with the underlying hardware, such as operating systems and device drivers. Built into the language is a documentation generator called Ddoc. Memory management. Memory is usually managed with garbage collection, but specific objects can be finalized immediately when they go out of scope. Explicit memory management is possible using the overloaded operators new and delete, as well as simply calling C's malloc and free directly. It is also possible to disable garbage collection for individual objects, or even for the entire program if more control over memory management is desired. Interaction with other systems. C's ABI (Application Binary Interface) is supported as well as all of C's fundamental and derived types, enabling direct access to existing C code and libraries. C's standard library is part of standard D. In D 1.0, C++'s ABI is not supported, although it can access C++ code that is written to the C ABI, and can access C++ COM (Component Object Model) code. D 2 already supports some interaction with the C++ ABI. Implementation. Current D implementations compile directly into native code for efficient execution. Getting and installing D. The Digital Mars D compiler can be obtained from the digital mars website. http://www.digitalmars.com/d/download.html You will need the two files dmd.zip and dmc.zip. According to the manual both files should be extracted into a root directory or one without spaces or other special characters. The location of link.exe should then be added to the path. D programs can now be compiled by calling 'dmd'. Win32: Example of configuration 1 (dmd). Create a batch file "dmd_vars.bat" and move it to a directory that is included in the path: @echo off echo Setting up a dmd environment... set PATH=c:\dm\bin;c:\dmd\bin rem ;%SystemRoot%\System32 set LIB=c:\dmd\lib;c:\dm\lib echo PATH set to %PATH% Then: Note: This works well, even if you have other compilers/tools installed (which might also have link.exe/make.exe etc.) Your first D programs. /First Program Examples/ With Tango library the classic hello world program is:  import tango.io.Console;  void main()  Cout("Hello, World").newline;  /Alternate Program Examples/ With the Phobos library the classic hello world program is:  import std.stdio;  void main()  writefln("Hello, World");  Compilation. Compiling hello world: dmd. dmd hello.d -ofhello gdc. gdc hello.d -o hello All D features. This is an incomplete list of all D features. It is specifically created to show and teach the D programming language with lots of examples. 

« Ethane | Stereoisomers » Many more linear alkanes can be formed by adding one additional carbon to the end of a chain of carbons. Ethane is the shortest chain with two carbons but DNA is known to have carbon chains containing millions of linked carbons. Alkanes are named based on the number of carbons in their longest chain: Naming carbon chains up to ten. The general equation for Alkanes is CnH2n+2.(Substute the n with no.of carbons, eg:-Propane 3 C 'C3H2*3+2=C3H8') The suffixis on the first four are from an obscure system but you should be familiar with the rest. Propane and butane are gases at standard temperature and pressure and are used commonly in lighters. Pentane on down the list are liquids at STP. Octane is the same as the octane in your gas tank. You may also want to name hydrocarbon chains of more than ten carbons. Here is a list for your reference: Naming longer carbon chains. etc... As the carbon chains get longer, molecules get relatively heavier and tend to move from being gases at Standard Temperature and Pressure to liquids to waxy solids. Foundational concepts » | « Alkanes | « Ethane | Propane thru decane | Drawing alkanes » | Introduction to reactions » 

Alkanes are the simplest organic molecules, consisting solely of singly-bonded carbon and hydrogen atoms. Alkanes are used as the basis for naming the majority of organic compounds (their nomenclature). Alkanes have the general formula CnH2n+2. Although their reactivities are often rather uninteresting, they provide an excellent basis for understanding bonding, conformation, and other important concepts which can be generalized to more "useful" molecules. =Introduction= Alkanes are the simplest and the least reactive hydrocarbon species containing only carbons and hydrogens. They are commercially very important, for being the principal constituent of gasoline and lubricating oils and are extensively employed in organic chemistry; though the role of pure alkanes (such as hexanes) is delegated mostly to solvents. The distinguishing feature of an alkane, making it distinct from other compounds that also exclusively contain carbon and hydrogen, is its lack of unsaturation. That is to say, it contains no double or triple bonds, which are highly reactive in organic chemistry. Though not totally devoid of reactivity, their lack of reactivity under most laboratory conditions makes them a relatively uninteresting, though very important component of organic chemistry. As you will learn about later, the energy confined within the carbon-carbon bond and the carbon-hydrogen bond is quite high and their rapid oxidation produces a large amount of heat, typically in the form of fire. As said it is important, not considered very important component in the chemistry. Introductory Definitions. Organic compounds contain carbon and hydrogen by definition and usually other elements (e.g. nitrogen and oxygen) as well. (CO2 is not an organic compound because it has no hydrogen). Hydrocarbons are organic compounds that contain carbon and hydrogen only. Alkanes are hydrocarbons or organic compounds made up of only carbon-carbon single bonds.Hence they are saturated. (as opposed to double and triple bonds). The simplest alkane is methane. Methane. Methane, (CH4, one carbon bonded to four hydrogens) is the simplest organic molecule. It is a gas at standard temperature and pressure (STP). This is a flattened, two-dimensional representation of methane that you will see commonly. The true three-dimensional form of methane does not have any 90 degree angles between bonded hydrogens. The bonds point to the four corners of a tetrahedron, forming cos-1(-1/3) ≈ 109.5 degree bond angles. Ethane. Two carbons singly bonded to each other with six hydrogens is called ethane. Ethane is the second simplest hydrocarbon molecule. It can be thought of as two methane molecules attached to each other, but with two fewer hydrogen atoms. Note that, if we were simply to create a new bond between the carbon centers of two methane molecules, this would violate the octet rule for the involved atoms. There are several common methods to draw organic molecules. They are often used interchangeably, although some methods work better for one situation or another. It is important to be familiar with the common methods, as these are the "languages" organic chemists can use to discuss structure with one another. =Drawing alkanes= When writing out the alkane structures, you can use different levels of the shorthand depending on the needs at hand in hand. For example, pentane can be written out. Its formula is C5H12. or CH3–CH2–CH2–CH2–CH3, or CH3(CH2)3CH3, or minimized to Line drawing shorthand. Although non-cyclic alkanes are called straight-chain alkanes they are technically made of linked chains. This is reflected in the line-drawing method. Each ending point and bend in the line represents one carbon atom and each short line represents one single carbon-carbon bond. Every carbon is assumed to be surrounded with a maximum number of hydrogen atoms unless shown otherwise. Structures drawn without explicitly showing all carbon atoms are often called "skeletal" structures, since they represent the skeleton or the backbone of the molecule. In organic chemistry, carbon is very frequently used, so chemists know that there is a carbon atom at the endpoints of every line that is not specifically labeled. =Conformations= Conformers, also called conformational isomers, or rotational isomers,or rotomers are arrangements of the same molecule made transiently different by the rotation in space about one or more single bonds. Other types of isomer can only be converted from one form to another by "breaking" bonds, but conformational isomers can be made simply by "rotating" bonds. Newman projections. Newman projections are drawings used to represent different positions of parts of molecules relative to each other in space. Remember that single bonds can rotate in space if not impeded. Newman projections represent different positions of rotating molecule parts. Conformations and energy. Different conformations have different potential energies. The staggered conformation is at a lower potential energy than the eclipsed conformation, and is favored. In ethane, the barrier to rotation is approximately 25 kJ/mol, indicating that each pair of eclipsed hydrogens raises the energy by about 8 kJ/mol. This number also applies to other organic compounds which have hydrogen atoms at similar distances from each other. At very low temperatures all conformations revert to the stabler( due to minimized vibration of atoms at it's mean position) , lower energy staggered conformation. Steric effects. Steric effects have to do with size. Two bulky objects run into each other and invade each others space. If we replace one or more hydrogen atoms on the above Newman projections with a methyl or other group, the potential energy goes up especially for the eclipsed conformations. Lets look at a Newman projection of butane as it rotates counterclockwise around its axes. When the larger groups overlap they repel each other more strongly than do hydrogen, and the potential energy goes up. Entropy. Entropy, represented as a ΔS, is a mathematical construct that represents disorder or probability. Natural systems want to find the lowest energy or organization possible, which translates to the highest entropy. "A note about potential energy: If you are rusty on this, remember the analogy of a big rock pushed to the top of a hill. At the top it has a maximum of potential energy. When you push it and allow it to roll down the hill the potential energy stored in it is transformed into kinetic energy that can be used to generate heat or smash something. " Notice that statistically, the ethane molecule has twice as many opportunities to be in the gauche conformation as in the anti conformation. However, because the Gauche configuration brings the methyl groups closer together in space, this generates high energy steric interactions and do not occur without the input of energy. Thus, the butane molecules shown will almost never be found in such unfavorable conformations. = Preparation of Alkanes = Wurtz reaction. Wurtz reaction is coupling of haloalkanes using sodium metal in solvent like dry ether  2R-X + 2Na → R-R + 2Na+X− Mechanism. The reaction consists of a halogen-metal exchange involving the free radical species R• (in a similar fashion to the formation of a Grignard reagent and then the carbon-carbon bond formation in a nucleophilic substitution reaction.) One electron from the metal is transferred to the halogen to produce a metal halide and an alkyl radical. The alkyl radical then accepts an electron from another metal atom to form an alkyl anion and the metal becomes cationic. This intermediate has been isolated in a several cases. The nucleophilic carbon of the alkyl anion then displaces the halide in an SN2 reaction, forming a new carbon-carbon covalent bond. Clemmensen reduction. Clemmensen reduction is a reduction of ketones (or aldehydes) to alkanes using zinc amalgam and hydrochloric acid The Clemmensen reduction is particularly effective at reducing aryl-alkyl ketones. With aliphatic or cyclic ketones, zinc metal reduction is much more effective The substrate must be stable in the strongly acidic conditions of the Clemmensen reduction. Acid sensitive substrates should be reacted in the Wolff-Kishner reduction, which utilizes strongly basic conditions; a further, milder method is the Mozingo reduction. As a result of Clemmensen Reduction, the carbon of the carbonyl group involved is converted from sp2 hybridisation to sp3 hybridisation. The oxygen atom is lost in the form of one molecule of water. Wolff-Kishner reduction. The Wolff–Kishner reduction is a chemical reaction that fully reduces a ketone (or aldehyde) to an alkane. Condensation of the carbonyl compound with hydrazine forms the hydrazone, and treatment with base induces the reduction of the carbon coupled with oxidation of the hydrazine to gaseous nitrogen, to yield the corresponding alkane. Mechanism. The mechanism first involves the formation of the hydrazone in a mechanism that is probably analogous to the formation of an imine. Successive deprotonations eventually result in the evolution of nitrogen. The mechanism can be justified by the evolution of nitrogen as the thermodynamic driving force. This reaction is also used to distinguish between aldehydes and ketones. Mozingo Reduction. A thioketal is first produced by reaction of the ketone with an appropriate thiol. The product is then hydrogenolyzed to the alkane, using Raney nickel = Properties of Alkanes = Alkanes are not very reactive when compared with other chemical species. This is because the backbone carbon atoms in alkanes have attained their octet of electrons through forming four covalent bonds (the maximum allowed number of bonds under the octet rule; which is why carbon's valence number is 4). These four bonds formed by carbon in alkanes are sigma bonds, which are more stable than other types of bond because of the greater overlap of carbon's atomic orbitals with neighboring atoms' atomic orbitals. To make alkanes react, the input of additional energy is needed; either through heat or radiation. Gasoline is a mixture of the alkanes and unlike many chemicals, can be stored for long periods and transported without problem. It is only when ignited that it has enough energy to continue reacting. This property makes it difficult for alkanes to be converted into other types of organic molecules. (There are only a few ways to do this). Alkanes are also less dense than water, as one can observe, oil, an alkane, floats on water. Alkanes are non-polar solvents. Since only C and H atoms are present, alkanes are nonpolar. Alkanes are immiscible in water but freely miscible in other non-polar solvents. Alkanes consisting of weak dipole dipole bonds can not break the strong hydrogen bond between water molecules hence it is not miscible in water. The same character is also shown by alkenes. Because alkanes contain only carbon and hydrogen, combustion produces compounds that contain only carbon, hydrogen, and/or oxygen. Like other hydrocarbons, combustion under most circumstances produces mainly carbon dioxide and water. However, alkanes require more heat to combust and do not release as much heat when they combust as other classes of hydrocarbons. Therefore, combustion of alkanes produces higher concentrations of organic compounds containing oxygen, such as aldehydes and ketones, when combusting at the same temperature as other hydrocarbons. The general formula for alkanes is CNH2N+2; the simplest possible alkane is therefore methane, CH4. The next simplest is ethane, C2H6; the series continues indefinitely. Each carbon atom in an alkane has sp³ hybridization. Alkanes are also known as paraffins, or collectively as the paraffin series. These terms are also used for alkanes whose carbon atoms form a single, unbranched chain. Branched-chain alkanes are called isoparaffins. Methane through Butane are very flammable gases at standard temperature and pressure (STP). Pentane is an extremely flammable liquid boiling at 36 °C and boiling points and melting points steadily increase from there; octadecane is the first alkane which is solid at room temperature. Longer alkanes are waxy solids; candle wax generally has between C20 and C25 chains. As chain length increases ultimately we reach polyethylene, which consists of carbon chains of indefinite length, which is generally a hard white solid. Chemical properties. Alkanes react only very poorly with ionic or other polar substances. The pKa values of all alkanes are above 50, and so they are practically inert to acids and bases. This inertness is the source of the term paraffins (Latin para + affinis, with the meaning here of "lacking affinity"). In crude oil the alkane molecules have remained chemically unchanged for millions of years. However redox reactions of alkanes, in particular with oxygen and the halogens, are possible as the carbon atoms are in a strongly reduced condition; in the case of methane, the lowest possible oxidation state for carbon (−4) is reached. Reaction with oxygen leads to combustion without any smoke; with halogens, substitution. In addition, alkanes have been shown to interact with, and bind to, certain transition metal complexes. Free radicals, molecules with unpaired electrons, play a large role in most reactions of alkanes, such as cracking and reformation where long-chain alkanes are converted into shorter-chain alkanes and straight-chain alkanes into branched-chain isomers. In highly branched alkanes and cycloalkanes, the bond angles may differ significantly from the optimal value (109.5°) in order to allow the different groups sufficient space. This causes a tension in the molecule, known as steric hinderance, and can substantially increase the reactivity. The same is preferred for alkenes too. =Introduction to Nomenclature= Before we can understand reactions in organic chemistry, we must begin with a basic knowledge of naming the compounds. The IUPAC nomenclature is a system on which most organic chemists have agreed to provide guidelines to allow them to learn from each others' works. Nomenclature, in other words, provides a foundation of language for organic chemistry. The names of all alkanes end with "-ane". Whether or not the carbons are linked together end-to-end in a ring (called "cyclic alkanes" or "cycloalkanes") or whether they contain side chains and branches, the name of every carbon-hydrogen chain that lacks any double bonds or functional groups will end with the suffix "-ane". Alkanes with unbranched carbon chains are simply named by the number of carbons in the chain. The first four members of the series (in terms of number of carbon atoms) are named as follows: Alkanes with five or more carbon atoms are named by adding the suffix "-ane" to the appropriate numerical multiplier, except the terminal "-a" is removed from the basic numerical term. Hence, C5H12 is called "pentane", C6H14 is called "hexane", C7H16 is called "heptane" and so forth. Straight-chain alkanes are sometimes indicated by the prefix "n-" (for normal) to distinguish them from branched-chain alkanes having the same number of carbon atoms. Although this is not strictly necessary, the usage is still common in cases where there is an important difference in properties between the straight-chain and branched-chain isomers: e.g. "n-hexane" is a neurotoxin while its branched-chain isomers are not. Number of hydrogens to carbons. This equation describes the relationship between the number of hydrogen and carbon atoms in alkanes: where "C" and "H" are used to represent the number of carbon and hydrogen atoms present in one molecule. If C = 2, then H = 6. Many textbooks put this in the following format: where "Cn" and "H2n+2" represent the number of carbon and hydrogen atoms present in one molecule. If Cn = 3, then H2n+2 = 2(3) + 2 = 8. (For this formula look to the "n" for the number, the "C" and the "H" letters themselves do not change.) Progressively longer hydrocarbon chains can be made and are named systematically, depending on the number of carbons in the longest chain. Naming carbon chains up to twelve. The prefixes of the first three are the contribution of a German Chemist, August Wilhelm Hoffman, who also suggested the name quartane for 4 carbons in 1866. However, the but- prefix had already been in use since the 1820s and the name quartane never caught on. He also recommended the endings to use the vowels, a, e, i (or y), o, and u, or -ane, -ene, -ine or -yne, -one, and -une. Again, only the first three caught on for single, double, and triple bonds and -one was already in use for ketones. Pent, hex, hept, oct, and dec all come from the ancient Greek numbers (penta, hex, hepta, octa, deka) and oddly, non, from the Latin novem. For longer-chained alkanes we use the special IUPAC multiplying affixes. For example, pentadecane signifies an alkane with 5+10 = 15 carbon atoms. For chains of length 30, 40, 50, and so on the basic prefix is added to -contane. For example, C57H116 is named as heptapentacontane. When the chain contains 20-29 atoms we have an exception. C20H42 is known as icosane, and then we have, e.g. tetracosane (eliding the "i" when necessary). For the length 100 we have "hecta" but for 200, 300 ... 900 we have "dicta", "tricta", and so on, eliding the "i" on "icta" when necessary; for 1000 we have "kilia" but for 2000 and so on, "dilia", "trilia", and so on, eliding the "i" on "ilia" when necessary. We then put all of the prefixes together in reverse order. The alkane with 9236 carbon atoms is then hexatridinoniliane. Isomerism. The atoms in alkanes with more than three carbon atoms can be arranged in many ways, leading to a large number of potential different configurations (isomers). So-called "normal" alkanes have a linear, unbranched configuration, but the "n-" isomer of any given alkane is only one of potentially hundreds or even possibly millions of configurations for that number of carbon and hydrogen atoms in some sort of chain arrangement.&lt;br&gt;Isomerism is defined as the compound having same moleculer formula the formula which present the different moleculer formula arrangement are called as Isomerism.&lt;br&gt;e.g.- The molecular formula for butane is C4H10. The number of isomers increases rapidly with the number of carbon atoms in a given alkane molecule; for alkanes with as few as 12 carbon atoms, there are over three hundred and fifty-five possible forms the molecule can take! Branched chains. Carbon is able to bond in all four directions and easily forms strong bonds with other carbon atoms. When one carbon is bonded to more than two other carbons it forms a branch. Above you see a carbon bonded to three and four other carbons. The common system has naming convention for carbon chains as they relate to branching. "Note: "R" in organic chemistry is a placeholder that can represent any carbon group." Constitutional isomers. One of the most important characteristics of carbon is its ability to form several relatively strong bonds per atom. It is for this reason that many scientists believe that carbon is the only element that could be the basis for the many complicated molecules needed to support a living being. One carbon atom can have attached to it not just the one or two other carbons needed to form a single chain but can bond to up to four other carbons. It is this ability to bond multiply that allows isomerism. Isomers are two molecules with the same molecular formula but different physical arrangements. Constitutional isomers have their atoms arranged in a different order. A constitutional isomer of butane has a main chain that is forked at the end and one carbon shorter in its main chain than butane. Naming Alkanes. There are several ways or systems for the nomenclature, or naming, of organic molecules, but just two main ones. The IUPAC system is necessary for complicated organic compounds. It gives a series of unified rules for naming a large compound by conceptually dividing it up into smaller, more manageable nameable units. Many traditional (non-IUPAC) names are still commonly used in industry, especially for simpler and more common chemicals, as the traditional names were already entrenched. IUPAC naming rules. Substituents are named like a parent, and replacing the "-ane" ending with "-yl". Numbering. The above molecule is numbered as follows: 2,3,7-Trimethyloctane Not 2,6,7-Trimethyloctane. Remember, number so as to give the smallest numbers to the substituents. Alphabetizing. 3-Ethyl-3-methylpentane "Ethyl" is listed before "methyl" for alphabetizing purposes. Branched Substituents. Naming branched substituents. 3-(1-methylethyl)-2,4-dimethylpentane The main chain in the drawing is numbered 1-5. The main part of the branched substituent, an ethyl group, is numbered 1' and 2'. The methyl substituent off of the ethyl substituent is not numbered in the drawing. To name the compound, put the whole branched substituent name in parentheses and then number and alphabetize as if a simple substituent. Common system. Some prefixes from the common system are accepted in the IUPAC system. For alphabetization purposes, iso- and neo- are considered part of the name, and alphabetized. Sec- and tert- are not considered an alphabetizable part of the name. Iso-. Iso- can be used for substituents that branch at the second-to-last carbon and end with two methyls. An isobutyl has four carbons total:  "Isobutyl" Sec-. Sec- can be used for substituents that branch at the first carbon . Neo-. Neo- refers to a substituent whose second-to-last carbon of the chain is trisubstituted (has three methyl groups attached to it). A neo-pentyl has five carbons total.  "Neopentyl" Tert-. Tert- is short for tertiary and refers to a substituent whose first carbon has three other carbon groups attached to it. =See also= 

Proton donors and acceptors » | Electrophiles and nucleophiles » The Lewis definition of acids and bases describes an acid as an electron acceptor and a base as an electron donator. « Foundational concepts | « Acids and bases | Proton donors and acceptors » | Electron donors and acceptors | Electrophiles and nucleophiles » | Alkanes » 

« Electron donors and acceptors | pKa and acidity » Electrophiles are "electron-lovers". (The suffix -phile means "lover of", as "bibliophile" means "lover of books"). Electrophiles seek electrons. Nucleophiles, or "nucleus lovers", seek positively charged nuclei. Electrophiles and nucleophiles are often ions, but sometimes not. « Foundational concepts | « Acids and bases | Electron donors and acceptors » | Electrophiles and nucleophiles | pKa and acidity » | Alkanes » 

« Electrophiles and nucleophiles | Alkanes » To measure acid strength of a compound, scientist typically use a quantity called pKa. It is defined as formula_1 where Ka is the acid dissociation constant. Stronger acids have small pKa values, where as weaker acids have higher pKa values. « Foundational concepts | « Acids and bases | Electrophiles and nucleophiles » | pKa and acidity | Alkanes » 

Stereoisomer of 2-bromo-3-hexanol &lt;/noinclude&gt; = Stereoisomers = Stereoisomers are a type of isomer where the order of the atoms in the two molecules is the same but their arrangement in space is different. To understand this we need to take a look at the ways that organic molecules can and cannot move. Again, usingthree-dimensional models is a great tool to visualize this and almost essential for most people to grasp these concepts clearly. With cyclo-alkanes, we observe that a group placed on one side of a ring stays on that same side. Except in very large rings (13+ carbons) the carbons are not free to rotate all of the way around their axes. This means that a group that is axial will not move into an equatorial position, and vice versa. Stereoisomerism is the arrangement of atoms in molecules whose connectivity remains the same but their arrangement in space is different in each isomer. The two main types of stereoisomerism are: Cis-trans Isomerism. "Main article:" Diastereomers Cis/trans isomerism occurs when a double bond is present, because the pi bond involved prevents that bond from being "twisted" the same way that a single bond can be. A good example is 1,2-dichloroethene: C2H2Cl2. Consider the two examples below: The two molecules shown above are "cis"-1,2-dichloroethene and "trans"-1,2-dichloroethene. This is more specifically an example of diastereomerism. These two molecules are stereoisomers because the two carbon atoms cannot be rotated relative to each other, due to the rigidity caused by the pi bond between them. Therefore, they are not "superimposeable" - they are not identical, and cannot take each other's place. However, the isomers are not mirror images of one another, so they are not enantiomers; therefore they must be diastereomers. Diastereomers usually have different chemical and physical properties and can exhibit dramatically different biological activity. There are two forms of these isomers; the "cis" and "trans" versions. The form in which the substituent hydrogen atoms are on the same side of the bond that doesn't allow rotation is called "cis"; the form in which the hydrogens are on opposite sides of the bond is called "trans". An example of a small hydrocarbon displaying cis-trans isomerism is 2-butene. Alicyclic compounds can also display cis-trans isomerism. As an example of a geometric isomer due to a ring structure, consider 1,2-dichlorocyclohexane: Optical Isomerism. "Main article:" Chirality Optical isomers are stereoisomers formed when asymmetric centers are present, for example, a carbon with four different groups bonded to it. Enantiomers are two optical isomers (i.e. isomers that are reflections of each other). Every stereocenter in one isomer has the opposite configuration in the other. Compounds that are enantiomers of each other have the same physical properties, except for the direction in which they rotate polarized light and how they interact with different optical isomers of other compounds. In nature, most biological compounds, such as amino acids, occur as single enantiomers. As a result, different enantiomers of a compound may have substantially different biological effects. When a molecule has more than one source of asymmetry, two optical isomers may be neither perfect reflections of each other nor superimposeable: some but not all stereocenters are inverted. These molecules are diastereomers, not enantiomers. Diastereomers seldom have the same physical properties. Optical isomerism is a form of isomerism (specifically stereoisomerism) where the two different isomers are the same in every way except being non-superposable [1] mirror images of each other. Optical isomers are known as chiral molecules. « Foundational concepts | « Alkanes | « Naming cycloalkanes | Stereoisomers | Introduction to reactions » &lt;/noinclude&gt; 

The links below will search on the named sites for the terms "organic chemistry model kit": 

Cycloalkanes are hydrocarbons containing one or more rings. (Alkanes without rings are referred to as aliphatic.)  → + H2 Under certain reaction conditions, propane can be transformed into cyclopropane. (H2 comes off as a sideproduct.)  Cyclopropane (unstable, lots of ring strain)  Cyclobutane (ring strain)  Cyclopentane (little ring strain)  Cyclohexane (Next to no ring strain)  Cyclodecane Rings with thirteen or more carbons have virtually no ring strain. Naming cycloalkanes. Cycloalkanes are named similarly to their straight-chain counterparts. Simply add the root "cyclo-" before the alkane part of the name. Example: Propane » Cyclopropane When naming cycloalkanes, the cyclo prefix is used for alphabetization. Substituents. If a cycloalkane has only one substituent, it is not necessary to assign that substituent a number. If there is more than one substituent, then it is necessary to number the carbons and specify which substituent is on which carbon.  Methylcyclopentane  1,1-dimethylcyclopentane  1,2-dimethylcyclopentane  1,3-dimethylcyclopentane The organic compound could be named and numbered  1-cyclopropyl-5-ethyl-2-methylcyclohexane but instead should be named  2-cyclopropyl-4-ethyl-1-methylcyclohexane because it produces a lower numbered name (1+5+2=8 vs. 2+4+1=7). In the following examples, notice that the longer chain is the parent and the cycloalkane is the substituent.  2-Cyclopropylbutane  1,3-dicyclopropylpropane Multicyclic alkanes. Multicyclic alkanes are hydrocarbons that have "more than one bonded cyclic ring". These abound in biology as all kinds of hormones, steroids, cholesterol,carbohydrates, etc. They are named as bicycloalkanes, tricycloalkanes, etc. They are named slightly differently than singularly cyclic alkanes.  Bicyclo[2.1.0]pentane Multicyclic alkanes are found frequently in living beings:  Part of Cholesterol We will get to some of the most interesting multicyclic rings later on when we study benzene and aromaticity. Stereochemistry. Because the C-C bonds in cycles cannot rotate through 360 degrees, substituted cycloalkanes and similar compounds can exhibit diastereomerism. This is comparable to alkenes which show cis/trans (or E/Z) isomerism. The isomers can be named using cis/trans notation, or more rigorously using R-S notation. Conformers, or conformational isomers, are different arrangements of the same molecule in space. Do not confuse them with any kind of true isomer as they are in every way the same molecule. The difference is in how the molecule is bent or twisted is space in any one instant of time. Cyclohexane. The first molecule that is generally presented in a discussion of cycloalkane conformers is cyclohexane. It comes in several flavors; the main ones are the chair conformation and the boat conformation. "Note: In the above models, the straight lines represents single bonds, the lumps represent carbon atoms, and the open ends represent hydrogen atoms." Consider getting a good molecular model set if you do not yet have one. They are not as inexpensive as you would hope but they help most people immensely to understand the way molecules look in three dimensions. Follow this link to places you can buy a molecular model kit. The chair conformation (can you see how it looks like a chair?) is lower in energy than the boat conformation. This is because the two ends of the molecule are farther apart and avoid steric hindrance. Hydrogen atoms in a cyclohexane can be divided into two types: When hydrogens are replaced with other, bulkier groups, it becomes apparent that the axial positions are less energetically favored than the equatorial positions. That means that, if given a choice, bulkier groups will tend to bond to cyclohexane in equatorial positions, as this reduces their steric hinderance and potential energy. Other cycloalkanes. Cyclopentane flips between slightly different conformers as well. 

Enantiomers. Main article: Chirality In chemistry, two stereoisomers are said to be enantiomers if they are mirror images of each other. Much as a left and right hand are different but one is the mirror image of the other, enantiomers are stereoisomers whose molecules are nonsuperposable mirror images of each other. 

Investing involves using your money (or borrowed money that you control) to earn more money. Before proceeding, make sure that you understand the concepts of Personal Finance. Topics covered here are: Many excellent resources available on the web: 

Some of the good news about organic chemistry is that it focuses on a subset of the periodic table, so there are fewer elements to worry about. These are the main elements with which we must concern ourselves, although occasionally a few others will be used in reactions. 

Limits, the first step into calculus, explain the complex nature of the subject. It is used to define the process of derivation and integration. It is also used in other circumstances to intuitively demonstrate the process of "approaching". Introduction. Intuitive Look into Limits. The limit is one of the greatest tools in the hands of any mathematician. We will give the limit an approach. Because mathematics came only due to approaches... remember?! We designate limit in the form: This is read as "The limit of formula_2 of formula_3 as formula_3 approaches formula_5". This is an important thing to remember, it is basic notation which is accepted by the world. We'll take up later the question of how we can determine whether a limit exists for formula_6 at formula_5 and, if so, what it is. For now, we'll look at it from an intuitive standpoint. Let's say that the function that we're interested in is formula_8 , and that we're interested in its limit as formula_3 approaches formula_10. Using the above notation, we can write the limit that we're interested in as follows: One way to try to evaluate what this limit is would be to choose values near 2, compute formula_6 for each, and see what happens as they get closer to 2. There are two ways to approach values near 2. One is to approach from below, and the other is to approach from above: The table above is the case from below. The table above is the case from above. We can see from the tables that as formula_3 grows closer and closer to 2, formula_6 seems to get closer and closer to 4, regardless of whether formula_3 approaches 2 from above or from below. For this reason, we feel reasonably confident that the limit of formula_16 as formula_3 approaches 2 is 4, or, written in limit notation, We could have also just substituted 2 into formula_16 and evaluated: formula_20. However, this will not work with all limits. Now let's look at another example. Suppose we're interested in the behavior of the function formula_21 as formula_3 approaches 2. Here's the limit in limit notation: Just as before, we can compute function values as formula_3 approaches 2 from below and from above. Here's a table, approaching from below: And here from above: In this case, the function doesn't seem to be approaching a single value as formula_3 approaches 2, but instead becomes an extremely large positive or negative number (depending on the direction of approach). Well, one says such a limit does not exist because no finite number is approached. This arises the concept of infinity: an undefined quantity and the limit is also called infinite limit or limit without a bound. Note that we cannot just substitute 2 into formula_26 and evaluate as we could with the first example, since we would be dividing by 0. Both of these examples may seem trivial, but consider the following function: This function is the same as Note that these functions are really completely identical; not just "almost the same," but actually, in terms of the definition of a function, completely the same; they give exactly the same output for every input. In elementary algebra, a typical approach is to simply say that we can cancel the term formula_29 , and then we have the function formula_8. However, that would be inaccurate; the function that we have now is not really the same as the one we started with, because it is defined when formula_31 , and our original function was specifically not defined when formula_31. This may seem like a minor point, but from making this kind of assumptions we can easily derive absurd results, such that formula_33 (see Mathematical Fallacy § All numbers equal all other numbers in Wikipedia for a complete example). Even without calculus we can avoid this error by stating that: In calculus, we can introduce a more intuitive and also correct way of looking at this type of function. What we want is to be able to say that, although the function isn't defined when formula_31, it works almost as if it was. It may not get there, but it gets really, really close. For instance, formula_36. The only question that we have is: what do we mean by "close"? Informal Definition of a Limit. As the precise definition of a limit is a bit technical, it is easier to start with an informal definition; we'll explain the formal definition later. We suppose that a function formula_2 is defined for formula_3 near formula_39 (but we do not require that it be defined when formula_40). Notice that the definition of a limit is not concerned with the value of formula_6 when formula_40 (which may exist or may not). All we care about are the values of formula_6 when formula_3 is close to formula_39 , on either the left or the right (i.e. less or greater). Limit can also be understood as: formula_3 is infinitely approaching to formula_39 but never equals to formula_39, just like the function formula_49, which infinitely approaches to formula_50 but never equals formula_50. Basics. Rules and Identities. Now that we have defined, informally, what a limit is, we will list some rules that are useful for working with and computing limits. You will be able to prove all these once we formally define the fundamental concept of the limit of a function. First, the constant rule states that if formula_52 (that is, formula_2 is constant for all formula_3) then the limit as formula_3 approaches formula_39 must be equal to formula_57. In other words Second, the identity rule states that if formula_59 (that is, formula_2 just gives back whatever number you put in) then the limit of formula_2 as formula_3 approaches formula_39 is equal to formula_39. That is, The next few rules tell us how, given the values of some limits, to compute others. Notice that in the last rule we need to require that formula_66 is not equal to 0 (otherwise we would be dividing by zero which is an undefined operation). These rules are known as identities; they are the scalar product, sum, difference, product, and quotient rules for limits. (A scalar is a constant, and, when you multiply a function by a constant, we say that you are performing scalar multiplication.) Using these rules we can deduce another. Namely, using the rule for products many times we get that This is called the power rule. As a result, we can safely say that all limits for polynomial functions can be deduced into several limits that satisfy the identity rule and thus easier to compute. Find the limit formula_69. We need to simplify the problem, since we have no rules about this expression by itself. We know from the identity rule above that formula_70. By the power rule, formula_71. Lastly, by the scalar multiplication rule, we get formula_72. formula_73 Find the limit formula_74. To do this informally, we split up the expression, once again, into its components. As above, formula_75. Also formula_76 and formula_77. Adding these together gives Find the limit formula_80. From the previous example the limit of the numerator is formula_81. The limit of the denominator is As the limit of the denominator is not equal to zero we can divide. This gives Find the limit formula_85. We apply the same process here as we did in the previous set of examples; We can evaluate each of these; formula_87 Thus, the answer is formula_88. Find the limit formula_90. In this example, evaluating the result directly will result in a division by 0. While you can determine the answer experimentally, a mathematical solution is possible as well. First, the numerator is a polynomial that may be factored: formula_91 Now, you can divide both the numerator and denominator by formula_29: formula_93 Remember that the limit is a method to determine the approaching value of a function instead of the value of the function itself. So, though the function is undefined at formula_31, the value of the function is approaching to formula_95 when formula_96 Find the limit formula_98. To evaluate this seemingly complex limit, we will need to recall some sine and cosine identities (see Chapter ). We will also have to use two new facts. First, if formula_6 is a trigonometric function (that is, one of sine, cosine, tangent, cotangent, secant and cosecant functions), and is defined at formula_5 , then formula_101. Second, formula_102. This can be proved using squeeze theorem. Note that L'Hospital's rule is not allowed to be used to evaluate this limit because it causes circular reasoning, in the sense that differentiating formula_103.requires this limit to equal one, which is exactly the equation that is being proven. Method 1 (before learning L'Hôpital's rule): To evaluate the limit, recognize that formula_104 can be multiplied by formula_105 to obtain formula_106 which, by our trig identities, is formula_107. So, multiply the top and bottom by formula_105. (This is allowed because it is identical to multiplying by one.) This is a standard trick for evaluating limits of fractions; multiply the numerator and the denominator by a carefully chosen expression which will make the expression simplify somehow. In this case, we should end up with: Our next step should be to break this up into formula_109 by the product rule. As mentioned above, formula_110. Next, formula_111. Thus, by multiplying these two results, we obtain 0. formula_73 Note that we also cannot apply L'Hospital's rule to evaluate this limit because it causes circular reasoning. We will now present an amazingly useful result, even though we cannot prove it yet. We can find the limit at formula_39 of any polynomial or rational function, as long as that rational function is defined at formula_39 (so we are not dividing by 0). That is, formula_39 must be in the domain of the function. We already learned this for trigonometric functions, so we see that it is easy to find limits of polynomial, rational or trigonometric functions wherever they are defined. In fact, this is true even for combinations of these functions; thus, for example, formula_116. The Squeeze Theorem. The Squeeze Theorem is very important in calculus, where it is typically used to find the limit of a function by comparison with two other functions whose limits are known. It is called the Squeeze Theorem because it refers to a function formula_2 whose values are squeezed between the values of two other functions formula_118 and formula_119 , both of which have the same limit formula_120. If the value of formula_2 is trapped between the values of the two functions formula_118 and formula_119 , the values of formula_2 must also approach formula_120. Expressed more precisely: Example: Compute formula_126. Since we know thatformula_127Multiplying formula_3 into the inequality yieldsformula_129Now we apply the squeeze theoremformula_130Since formula_131, formula_132 formula_73 Finding Limits. Now, we will discuss how, in practice, to find limits. First, if the function can be built out of rational, trigonometric, logarithmic, or exponential functions, then if a number formula_39 is in the domain of the function, then the limit at formula_39 is simply the value of the function at formula_39:formula_137 when formula_6 can be built out of rational, trigonometric, logarithmic, or exponential functions and formula_139 the Domain of formula_6If formula_39 is not in the domain of the function, then in many cases (as with rational functions) the domain of the function includes all the points near formula_39, but not formula_39 itself. An example would be if we wanted to find formula_144 , where the domain includes all numbers besides 0. In that case, in order to find formula_145 we want to find a function formula_146 similar to formula_6 , except with the hole at formula_39 filled in. The limits of formula_2 and formula_118 will be the same, as can be seen from the definition of a limit. By definition, the limit depends on formula_6 only at the points where formula_3 is close to formula_39 but not equal to it, so the limit at formula_39 does not depend on the value of the function at formula_39. Therefore, if formula_156 , formula_157 also. And since the domain of our new function formula_118 includes formula_39 , we can now (assuming formula_118 is still built out of rational, trigonometric, logarithmic and exponential functions) just evaluate it at formula_39 as before. Thus we have formula_162. In our example, this is easy; canceling the formula_3's gives formula_164, which equals formula_165 at all points except 0. Thus, we have formula_166. In general, when computing limits of rational functions, it's a good idea to look for common factors in the numerator and denominator. Specific DNE Situations. Note that the limit might not exist at all (DNE means "does not exist"). There are a number of ways in which this can occur: "Gap" Determining Limits. The formal way to determine whether a limit exists is to find out whether the value of the limit is the same when approaching from below and above (see at the top of this chapter). The notation for the limit approaching from below (in increasing order) isformula_209, notice the negative sign on the constant formula_39The notation for the limit approaching from above (from decreasing order) isformula_211, notice the positive sign on the constant formula_39For example, let us find the limits of formula_49 when formula_3 is approaching formula_50 in both directions. In other words, find formula_216 and formula_217.Recall the table we made when we are trying to intuitively feel the limit. We can use that to help us. However, if familiar enough with reciprocal functions, we can simply determine the two values by imagining the graph of the function. The following table is when formula_3 is approaching from below. Thus, we found that when formula_3 is approaching from below to formula_50, the function approaches negative infinity. In mathematical terms: formula_221 Now let's talk about the approach from above. We found that formula_222formula_73 The method of determining if limits exist is relatively intuitive. It only requires some practice to be familiar with the process. Let's use the same example: find formula_224.Since we already found that formula_221 and formula_222, and obviously, formula_227 We can say that formula_224 does not exist.formula_73 Infinity Situations. Now consider the function What is the limit as formula_3 approaches zero? The value of formula_232 does not exist; it is not defined. Notice, also, that we can make formula_146 as large as we like, by choosing a small formula_3 , as long as formula_235. For example, to make formula_146 equal to formula_237 , we choose formula_3 to be formula_239. Thus, formula_240 does not exist. However, we "do" know something about what happens to formula_146 when formula_3 gets close to 0 without reaching it. We want to say we can make formula_146 arbitrarily large (as large as we like) by taking formula_3 to be sufficiently close to 0, but not equal to 0. We express this symbolically as follows: Note that the limit does not exist at formula_50 ; for a limit, being formula_247 is a special kind of not existing. In general, we make the following definition. An example of the second half of the definition would be that formula_248. Applications of Limits. To see the power of the concept of the limit, let's consider a moving car. Suppose we have a car whose position is linear with respect to time (that is, a graph plotting the position with respect to time will show a straight line). We want to find the velocity. This is easy to do from algebra; we just take the slope, and that's our velocity. But unfortunately, things in the real world don't always travel in nice straight lines. Cars speed up, slow down, and generally behave in ways that make it difficult to calculate their velocities. Now what we really want to do is to find the velocity at a given moment (the instantaneous velocity). The trouble is that in order to find the velocity we need two points, while at any given time, we only have one point. We can, of course, always find the average speed of the car, given two points in time, but we want to find the speed of the car at one precise moment. This is the basic trick of differential calculus, the first of the two main subjects of this book. We take the average speed at two moments in time, and then make those two moments in time closer and closer together. We then see what the limit of the slope is as these two moments in time are closer and closer, and say that this limit is the slope at a single instant. We will study this process in much greater depth later in the book. First, however, we will need to study limits more carefully. 

Welcome to the Wikibook about Artificial Intelligence. Book Contents. The following is a first proposal for a basic layout. This is not yet complete, ideas are welcome. Discuss on the talk page or just add them here. The book is laid out into 5 sections, with increasing detail and complexity. Each section contains a number of chapters. In addition to regular chapters, there are case-study chapters that investigate full and complex AI systems using several techniques from the regular chapters (as well as perhaps some new ones). Introduction. Overview 

¡Bienvenidos! "Welcome!" Vaya a los contenidos» "Go to the contents»" 

Introduction. This is the very first lesson in learning a second language, the Spanish language! This lesson begins with simple greetings, and covers important ideas of the Spanish language. Throughout education, methods of teaching Spanish have changed greatly. Years ago, the Spanish language was taught simply by memory. Today, however, the Spanish Language is taught by moving more slowly and covering grammar and spelling rules. Again, this is an introduction. If this is the first time you are attempting to learn Spanish, do not become discouraged if you cannot understand, pronounce, or memorize some of the things discussed here. In addition, learning a second language requires a basic understanding of your own language. You may find, as you study Spanish, that you learn a lot about English as well. At their core, all languages share some simple components like verbs, nouns, adjectives, and plurals. English, as your first language, comes naturally to you and you don't think about things like subject-verb agreement, verb conjugation, or usage of the various tenses; yet you use these concepts on a daily basis. While English is described as a very complicated language to learn, many of the distinguishing grammar structures have been simplified over the years. This is not true for many other languages. Following the grammatical conventions of Spanish will be very important, and can actually change the meaning of phrases. You'll see what is meant by this as you learn your first verbs ser and estar. Do not become discouraged! You can do it. Dialogue 1. Two good friends - Carmen and Roberto - are meeting:  (139KB) Dialogue 2. Two people - Señor González and Señora Pérez - are meeting for the first time: Nice to meet you. Listen to the Dialogue. Vocabulary. Exercise: Greetings Grammar: Subject Pronouns. &lt;/br&gt; A few things to keep in mind: Grammar: Verbs ser and estar. Spanish has two different words that can be translated with ""to be". Ser is used more for more permanent characteristics ("Soy Luis") whereas estar is used for more temporary or changeable conditions, such as location ("La papelera está al lado del escritorio"", "The trash can is beside the desk") and feeling (""Estoy bien""). A good way to remember when to use "estar" is by using the rhyme, "To tell how you feel or where you are, always use the verb estar." In future lessons we will come back to the uses of ser and estar. Here we will look at the conjugations in the "present indicative". &lt;br&gt; Ejemplos de los verbos ser y estar (Examples of the verbs ser and estar). Note: *This use of estar is the Spanish "present progressive" which is used for actions in progress. More about the "present progressive" in Lesson 4 Dialect Note: Spanish which uses the "vos" form conjugates ser with the following irregular form: "sos". Exercise: Verbs ser and estar Hay. Spanish uses a different verb (haber) to express ""there is " and "there are". The form of haber used for this purpose is hay, for both singular ("there is") and plural ("there are"). Spanish alphabet. Here is the normal Spanish alphabet. However words aren't alphabetized by it. Please read the notes and sections below. (Blue letters are a part of the normal English alphabet.) Audio: (646KB) Although the above will help you understand, proper pronunciation of Spanish consonants is a bit more complicated: Most of the consonants are pronounced as they are in American English with these exceptions: Vowel pronunciation. The pronunciation of vowels is as follows: The "u" is always silent after q (as in "qué" pronounced kā). Spanish also uses the ¨ (diaeresis) diacritic mark over the vowel u to indicate that it is pronounced separately in places where it would normally be silent. For example, in words such as "vergüenza" ("shame") or "pingüino" ("penguin"), the "u" (sounds the same as the "u" in "ultra") is pronounced similarly but with more strength to the English "w" forming a diphthong with the following vowel: [we] and [wi] respectively. It is also used to preserve sound in stem changes and in commands. Semi-Vowels. At the end of a word or when it means "and" ("y") it is pronounced like i. Acute accents. Spanish uses the ´ (Acute) diacritic mark over vowels to indicate a vocal stress on a word that would normally be stressed on another syllable; Stress is contrastive. For example, the word "ánimo" is normally stressed on "a", meaning "mood, spirit." While "animo" is stressed on "ni" meaning "I cheer." And "animó" is stressed on "mó" meaning "he cheered." Additionally the acute mark is used to disambiguate certain words which would otherwise be homographs. It's used in various question word or relative pronoun pairs such as "cómo" (how?)&amp; "como" (as), "dónde"(where?) &amp; "donde" (where), and some other words such as "tú" (you) &amp; "tu" (your), "él" (he/him) &amp; "el" (the). Emphasis. The rules of stress in Spanish are: 1. When the word ends in a vowel or in "n" or "s" the emphasis falls on the second to last syllable. Eg: Mañana, Como, Dedos, Hablan. 2. When the word ends in a consonant other than "n" or "s", the emphasis falls on the last syllable. Eg: Ciudad, Comer, Reptil. 3. If the above two rules don't apply, there will be an accent to show the stress. Eg: Fíjate, Inglés, Teléfono. 4. SPECIAL CASE: Adverbs ending in "-mente", which are derived from adjectives, have two stresses. The first stress occurs in the "adjective part" of the adverb, on the syllable where the adjective would normally be stressed. The second stress occurs on the "-men-" syllable. Eg: Solamente, Felizmente, Cortésmente. 

Díalogo 1: El salón de clase. "La profesora entra en el salón." "Profesora": Buenos días alumnos. Hoy estudiamos los objetos en el salón de clase. ¿Cuántos libros hay? "Carlos": Hay catorce estudiantes. Hay catorce libros. "Profesora": Bien Carlos. Marianela, ¿Hay un mapa en el salón? "Marianela": Sí, señora. Hay un mapa a la izquierda de la pizarra. Está en la pared. "Profesora": Bien, Marianela. ¿De qué color es el mapa? "Carlos": El mapa es verde y azul. Es el mapa de México. "Marianela": ¿Profesora, porqué está el mapa en la pared y no está en la pizarra? "Profesora": La pizarra es para escribir con la tiza. Tengo un borrador aquí. ¿Quieres escribir la capital de México en la pizarra? "Marianela": No sé la capital de México. "Carlos": Marianela, mira el mapa. "Marianela": Ay, gracias Carlos. La capital es Ciudad de México. "Profesora": Bien. Vocabulario: El salón de clase. All vocabulary will be given with the appropriate definite article; however, the article will not be translated. In addition, note that all nouns have a gender: they are either masculine or feminine. This will be important in later lessons when you begin to learn adjectives. Grammar: The Definite Article. Like in English, Spanish has definite articles that serve to identify the location of the noun in the sentence. These are commonly called "noun markers". These articles have a gender that is equivalent to the gender of the noun they modify. However, sometimes pronunciation requires the article to be different than the gender of the word. In these situations, the gender of the noun will always be indicated. Otherwise, you can use the article to determine the gender of the noun. &lt;br&gt; Examples: The following examples come from the vocabulary of the first dialogue. Note that no matter whether the word is singular or plural, the meaning of the definite article does not change. &lt;br&gt; The Indefinite Article. In English the indefinite articles are "a" and "an" (singular) or "some" (plural). In Spanish there are different forms for masculine-gender, feminine-gender, singular or plural. &lt;br&gt; Examples: &lt;br&gt; &lt;br&gt; For phonetic reasons some words beginning with accented a may have the article un: un ave blanca "(a white bird)," las aves blancas "(the white birds)". This is basically the same idea as el ave blanca "(the white bird)". Remember, do not confuse uno "(one)" with un "(a or an)".&lt;br&gt; Also, do not confuse una "(a or an)" with uña "(fingernail)". Exercise: Spanish/Exercises/Articles The colors. As in English, colors in Spanish are adjectives. As adjectives, their endings will vary according to the nouns they modify. The following adjectives are in the masculine singular form. Most adjectives will be presented to you in this form. &lt;br&gt; Examples: &lt;br&gt; &lt;br&gt; Los números (Numbers). The numbers 16 through 29 can be formed in two ways: (1) using the conjunction "y" like veinte y uno - 21, or (2) as one complete word, changing the spelling of the word like dieciocho (18) or veintiséis (26). 





Numerical Axioms. It is possible to define a regular set of numbers in a formal fashion. The set of "Peano axioms" define the series of numbers known as the natural numbers. They are as follows: Let us attempt to motivate these axioms. We want these axioms to eliminate any set which is not the natural numbers. E.g., any set fulfilling the above should at least be infinite. The first two are obvious properties of the natural numbers (and of integers) as we know them. Note that some prefer to use 1 as the lowest number. The reason for choosing zero has root in [set theory], in which the first natural number is chosen as the empty set formula_1. The 3rd axiom prevents circularity. If this axiom was not included, defining formula_2 would trivially fulfill the remaining axioms --- prove this for yourself by considering each remaining axiom! The 4th prevents a partial loop. Consider a the set formula_3 and set formula_4 and formula_5. This set fulfills every axiom but the 4th --- prove this for yourself. The 5th is sometimes called the induction axiom. It ensures that the set is "connected", i.e. that we can reach any number by using the 2nd axiom repeatedly on 0. An example of a set that fulfills every axiom but the 5th is formula_6 with the usual meaning of +1. From this we can deduce the existence of a series of quantities like this: where '0' is a constant and the first natural number and '1' is a constant natural number equivalent to the difference in value between two consecutive natural numbers. This set is sufficient for counting. However, it is inconvenient to refer to a large natural number as '0' followed by the requisite large number of '+ 1' expressions. Due to this, each of the natural numbers is given a label, and to make the labelling easier another axiom is introduced: '0 + 1' is equivalent to '1'. Thus the series of natural numbers may be written so for some brevity: Once this is done, giving each quantity its own label is trivial. And so the series of natural numbers can then be written: 

formula_1 formula_2 Discrete mathematics is the study of mathematical structures that are fundamentally discrete rather than continuous. Contents.  __NOEDITSECTION__ 

The subject of mathematics is commited to rigorous reasoning. This book aims to raise your confidence in the manipulation and interpretation of formal notations, as well as to train you to solve problems that are fundamentally discrete: problems like chess, in which the moves you make are exact; problems where fields like calculus fail because there's no continuity; problems that appear frequently in games, puzzles, and computer science. We hope you'll enjoy discovering discrete mathematics here, and we hope you'll find this a good reference for quickly picking up the details you may forget. 

"This is currently a pre-draft textbook. Therefore some information is inaccurate." Coi! This book will help you learn Lojban. Old Table Of Contents. Appendices 

Lojban has a 24-letter alphabet (23 of which come from the Roman alphabet):&lt;br&gt; &lt;br&gt; ' a b c d e f g i j k l m n o p r s t u v x y z&lt;br&gt; &lt;br&gt; And two special symbols:&lt;br&gt; &lt;br&gt; , .&lt;br&gt; &lt;br&gt; This includes most of the traditional English alphabet except for three letters:&lt;br&gt; &lt;br&gt; h q w&lt;br&gt; &lt;br&gt; Lojban has a phonetic spelling and audio-visual isomorphism, meaning that each letter has exactly one sound and each sound has exactly one letter.&lt;br&gt; &lt;br&gt; ' only appears in Lojban words between vowels. "," only appears in Lojban names. Lojbanists can optionally not include the "." in their writing. The stress is on a word's penultimate syllable. The exception is that a syllable whose nucleus is a 'y' cannot be stressed; in that case, stress the syllable immediately preceding it. Names may indicate different from penultimate stress by capitalizing the entire stressed syllable or the vowel of the stressed syllable. Some Lojbanists have unofficially proposed other characters to mark stress like "^".&lt;br&gt; &lt;br&gt; Lojban has eleven diphthongs written by composed vowels. The following transliterations are assuming that the reader has an American English accent.&lt;br&gt; A pair of adjacent vowels which are capable of forming a diphthong "will" form a diphthong unless a comma is placed between them. That is the purpose of the comma in Lojban; to act as a separator between syllables whenever necessary for pronunciation purposes. (It is not for pausing between words or phrases like in English.) To make pronouncing some words easier, you may put unwritten buffer vowels in your speech between consonants. Any non-Lojban vowel can act as a buffer vowel, but i as in h"i"t has the most popularity. For example, you could pronounce (but not write) a buffer vowel between the "s" and "f" of "sfani".&lt;br&gt; &lt;br&gt; Exercises. Write/Stress Answer True or False Answers to exercises. Pronounce Write/Stress Answer True or False 

In Lojban names are called "cmene" (pronouced shmen-eh, IPA /ʃme.ne/), and are transliterated to use only Lojban sounds. In Lojban all words end in a vowel, and so to avoid any possible confusion all cmene end in a consonant; if a transliterated name would end in a vowel, it is usual to add an s to the end. All cmene are followed by a full stop, so confusion to where the name stops and the next word starts it avoided; for this reason, any names beginning in a vowel have a full stop in front of them. To indicate an accent somewhere other than the second-to-last syllable, the entire syllable to be accented must be capitalised. A comma should be used to stop two vowels next to each other that should not be pronouced as a dipthong. All names must have the word 'la', meaning 'the', placed before them. A few examples are below: Exercises. English names to Cmene. Translate these names to Lojban cmene. Cmene to English names. Translate these Lojban cmene to English names. Answers to Exercises. No one pronouciation of a name is the right one, so use common sense whe deciding if you got these right or not. 

Stem Changing Verbs. In Spanish, some verbs change their stems when they are conjugated. These verbs are known as "stem-changing verbs". Many of these verbs are important and often used. There are three different types of stem changing verbs in Spanish: The stem changes for all conjugations, excepting nosotros/as and vosotros/as. The endings are the same as for regular verbs (-o for yo, -as/-es for tú, ...). &lt;br&gt; Entender&lt;br&gt;To Understand Note that the stem change is done for the second 'e' (not the first one) - in general the stem always changes for the last vowel before the -ar/-er/-ir ending. Example: pedir (e-&gt;i) "to ask for, to order"  Pedir&lt;br&gt;To Ask For, Order Note: all e-&gt;i stem changing verbs are -ir verbs.  Dormir&lt;br&gt;To Sleep Here is a list of some other common stem changing verbs: Exercise: Stem Changing Verbs Present Participle. The present participle in Spanish is used either for continuous tenses (with "estar", e.g. "I am running.") and can also be used as an adjective. The Spanish present participle corresponds to the English -ing form of the verb.To form the present participle for regular -ar verbs, add -ando to the stem. For -er and -ir verbs add -iendo: &lt;br&gt; However, not all present participles are that regular. Some verbs add a "y," or change the spelling, to adhere to Spanish orthography (spelling) rules. Here is a list of some common verbs that have an irregular present participle: Present Progressive. Like in English, the Spanish present progressive is used to describe an "action in progress". It is formed by conjugating the verb estar and then adding the present participle: 

Text. Here are a couple of sentences and short dialogs about people planning/doing leisure activities. Besides the new vocabulary you should also have a look at how the verbs are conjugated depending on the subject of the sentence. As you may see, each verb is bolded. These verbs are "conjugated", that is, changed by the person(s) to which they are referring. Notice that subject pronouns are "not" necessary. Regular Verbs. Spanish has three different types of regular verbs: -ar, -er, and -ir verbs. The subject pronoun is not necessary and in conversational Spanish it is only used for emphasis. For this lesson, we will omit it. One can still use pronouns, however. The conjugation pattern is the following: &lt;/br&gt; As one can see, the endings for each person are different. This is similar to other Romance languages, such as Portuguese and Italian (the notable exception is French). This is the reason why we may omit the pronouns while we speak. Remember that sometimes it is best to clarify whether él, ella, or usted is speaking, because they share the same form. However, the context of the rest of the sentence sometimes clarifies this. There are a few steps involved with conjugating a verb. Here are the steps involved: Notice that there are only two differences between the conjugations of -er and -ir verbs. The "nosotros" (4) and "vosotros" (5) forms are the only differences. Those forms have an "i" in the stem instead of an "e." Exercise: Regular Verbs "G" Verbs. The verb "hacer" means to do or to make. Hacer is irregular in the first person singular form "(I)" of the present tense only. The irregular form is hago. Hacer is one of the many verbs in Spanish which add a "g" in the first person singular of the verb. This is the present indicative of the verb hacer. Note that the verb hacer is translated as "to do" and "to make" when referring to activities. But it can also be used to talk about some weather conditions: But: When speaking about the weather using hacer, the Ud. form (third singular form) is always used. El vocabulario (Vocabulary) - Los días (Days). Audio: (157KB) Una fiesta. Una fiesta entre amigos. Nosotros bailamos y lo pasamos bien en el jardín de esta casa.&lt;br&gt; "A party among friends. We dance and enjoy ourselves in the patio (garden) of this house." 

The Sapir-Whorf Hypothesis (SWH) is a hypothesis in linguistics, stating that there are notable differences in thought patterns of speakers of different languages, and that the way people's brains function is strongly affected by their native languages. It's a very controversial theory, championed by linguist Edward Sapir and his student Benjamin Whorf. First discussed by Sapir in 1929, the hypothesis became popular in the 1950s following posthumous publication of Whorf's writings on the subject. In 1955, Dr. James Cooke Brown created the Loglan language (which led to the offshoot Lojban) in order to test the hypothesis. After vigorous attack from followers of Noam Chomsky in the following decades, the hypothesis is now only believed by linguists with a grain of salt; that thought processes are somewhat affected by language, but that differences aren't that notable. Central to the Sapir-Whorf hypothesis is the idea of linguistic relativity: distinctions of meaning between related terms in a language are often arbitrary and particular to that language. Sapir and Whorf took this one step further by arguing that a person's whole world view is determined by the vocabulary and syntax available in his or her language. The extreme ("Weltanschauung") version of this idea, that all mental function is constrained by language, can be disproved through personal experience: people in every language occasionally struggle to express their exact thoughts, feeling constrained by the language. It's common to say or write something, only to correct one's self or further clarify meaning, especially to someone being explained to. These show that ideas are not merely words, because one can imagine something without being able to express it in words. The opposite extreme — that language does not influence thought at all — is also widely considered to be false. For example, it has been shown in studies that people's discrimination of similar colors can be influenced by their vocabulary for distinguishing said colors. Another study showed that deaf children of hearing parents are more likely to fail on some cognitive tasks unrelated to hearing, while deaf children of deaf parents succeed, due to parents being able to more extensively communicate. Computer programmers who know different programming languages often see the same problem in completely different ways. The Neuro-Linguistic Programming (NLP) analysis of the problem is direct: most people do notable thinking by talking to themselves and by imagining images and other sensory phantasms. To the extent that people think by talking to themselves, they are limited by their vocabulary and the structure of their language and linguistic habits. (However it should also be noted that everyone have idiolects, mental language patterns individual to them.) John Grinder, a founder of NLP, was a linguistics professor who perhaps unconsciously combined the ideas of Chomsky with the Sapir-Whorf hypothesis. A seminal NLP insight came from a challenge he gave to his students: coin a neologism to describe an idea for which you have no words. Student Robert Dilts gave an example by coining a word for the way people stare into space when they are thinking, and for the different directions they stare. These new words enabled users to describe patterns in the ways people stare into space, which led directly to NLP results — as notable a validation of the weak hypothesis as one could ask. Lojban is thus designed to test the Sapir-Whorf hypothesis, by attempting to expand the speakers' minds and express thoughts as conveniently as possible to see if there's any notable effects in the speakers' thought patterns/worldview. The only complication with this is that the third factor — the speakers all wanting to learn Lojban, an obscure language — could skew the results somewhat, but the only way to fix that is to get more Lojban speakers. External links. "This article was originally taken from Wikipedia and is licensed under the GNU FDL. Update as needed." 

Loglan. In 1955, sociologist from Philippines Dr. James Cooke Brown began work on Loglan, a constructed language designed for linguistic research, particularly investigation of the Sapir-Whorf hypothesis, a theory stating that linguistic structures affect people's thoughts. The object was to make a language so powerfully expressive for logic and calculation that people learning it would become measurably smarter if the hypothesis was true. He intended it to be as culturally neutral as possible, logically and linguistically powerful, incorporating all known expressive features of any language (e.g., compounded location tenses), metaphysically parsimonious (e.g., you are not required to express any feature of reality, as you are in English time-tenses of verbs), and totally regular and unambiguous. He even used maximally stable phonemes. The formal grammar was disambiguated mechanically (at first). The language's grammar is based on predicate logic (the name Loglan is short for "logical language"), which also makes it particularly suitable for human-computer communication, an application that led Robert Heinlein to mention the language in his novel The Moon is a Harsh Mistress. Dr. Brown founded The Loglan Institute to develop the language and other applications of it. He always considered the language a research project that was never complete, so although he released many papers about its design he never "released" it to become a usable language. A group of his followers later formed The Logical Language Group to create the language Lojban along the same principles, but with the intention to make it freely available and encourage its use as a real language. This latter group has a small but active community of speakers. Loglan allows very strange speech, because one can say things that are simply not meaningful. For example, you can literally say that John, a person, is a short word. Or one can directly and precisely say any of the many possible meanings of the English phrase "a pretty little girls school." In natural languages, the ambiguity of the grammar hides these odd meanings; but in Loglan they are all available and add expressive freedom to the language. This feature is so pronounced that people fluent in Loglan say impossible things as a sort of joke — a type of humour simply not supported by the linguistic machinery of natural languages. The oddest, most difficult thing for a speaker of an Indo-European language like English is that Loglan has no distinction between nouns or verbs, objects, direct objects, indirect objects, possessive forms or tenses. There are only predicates, with places for variables: for example, "botso": X buys item Y from seller P for price Q. There are prefixes to reorder predicates; for example, price would be "botso" with a little word to make price the first variable. With different reorderings, one speaks of buyer, bought-thing, or seller. Tenses for time, location, actor, type of action, etc. are provided by "little words" which are optional. Predicates compound, so a predicate can fit in the variable of another predicate. Every feature of the language has standard, regular forms for acting in compounds. For example, time-travel tenses are available trivially in Loglan (I did X from time Y to P in time Q.) using compounding forms normally used for location tenses. After long use of Loglan, the world gains a sort of timeless, objectless, actorless flavor. Time words and location words fall away except when needed to make a point, usually with emotional emphasis. It is rather easy to avoid blame for responsibilities in Loglan, and scheduling may be creatively ambiguous because the tenses are optional. The language is designed so that the patterns of phonemes always parse into words. Thus, one cannot babble Loglan, because even when run together, the language is still parsable. Lojban. The constructed language Lojban (SAMPA ['loZban]) was created by the Logical Language Group in 1987 based on the earlier Loglan, with the intent to make the language more complete, usable, and freely available. The language itself shares many of the features and goals of Loglan; in particular: While the initial goal of the Loglan project was to investigate the Sapir-Whorf hypothesis, the active Lojban community has additional goals for the language, including: Lojban Grammar. All words in Lojban belong to one of three overall categories: "brivla", for both common nouns and verbs; "cmene", for proper nouns; "cmavo", for structural particles: articles, numerals, tense indicators and other such modifiers. The "cmavo" are further subdivided into "selma'o", which are closer to the notion of parts of speech (e.g., UI includes interjections and discursives). There is no distinct class of words for adjectives or adverbs, unlike most Indo-European languages. Most types of noun phrases are always preceded by an article; there are different articles to indicate whether it is being treated as an individual, mass, set, or typical element. "Brivla" do not inflect for tense, person, or number; tense may optionally be indicated by separate "cmavo", and number may optionally be indicated in various ways. As befits a logical language, there is a large assortment of conjunctions. Logical conjunctions take different forms depending on whether they connect "sumti" (the equivalent of noun phrases), "selbri" (phrases that can serve as verbs; all "brivla" are "selbri"), parts of a "tanru" (a construct whose closest English equivalent is a string of nouns), or clauses in a sentence. The typology is Subject Verb Object, with Subject Object Verb also common. Word formation is polysynthetic; many "brivla" (all of which, except for a handful of borrowings such as "alga", have at least five letters) have one to three three-letter forms called "rafsi" which are used in making longer words. For example, "gasnu" means "to make something happen"; its "rafsi" form of "-gau" regularly forms compounds meaning "to cause...x", in which the agent is in the subject place of the new predicate. Lojban has a positional case system, though this can be overridden by marking predicate arguments with explicit case particles. For instance, "bramau" means "is bigger than"; the bigger thing is in first position, and the smaller is second, and the measured property in the third. So "mi bramau do le ka clani" means "I am bigger than you in the property of height" or "I am taller than you"; but this could also be expressed as something like "fi le ka clani fe do fa mi bramau", "In height, you are exceeded by me". What a particular place means depends entirely on the "brivla". For animals and plants the second place is the species, variety, breed, or other taxon; for verbs of measurement it is the numerical measurement, and a further place is the standard; for "klama" ("go" / "come") it is the destination. There may be up to five places for some "brivla". Something of the flavor of Lojban (and Loglan) can be imparted by this lightbulb joke: Q: How many Lojbanists does it take to change a broken light bulb?&lt;br&gt; A: Two: one to decide what to change it into, and one to figure out what kind of bulb emits broken light. This makes use of two features of the language; first, the language attempts to eliminate polysemy, that is, having a word with more than one meaning. So while the English word "change" can mean "to transform into a different state", or "to replace", or even "small-denomination currency", Lojban has different words for each. In particular, the use of a "brivla" such as the word for "change" ("binxo") implies that all of its predicate places exist, so there must be something for it to change into. Another feature of the language is that it has no grammatical ambiguities such as appear in English phrases like "big dog catcher", which can mean either a big person who catches dogs or a person who catches big dogs. In Lojban, unless you clearly specify otherwise with "cmavo", such modifiers always group left-to-right, so "big dog catcher" is a catcher of big dogs, and a "broken light bulb" is a bulb that emits broken light (you can also avoid the ambiguity by creating a new word, so "broken lightbulb" has the intended meaning). 



Attitudinals are markers of attitude or emotion, modifying the word (brivla, cmavo, cmene, sumti, or tanru) directly before it. For example, the simple « "mlatu" » for "cat" can become « "mlatu .ui" » for "cat (happiness!/:D)". At the beginning of a sentence (after the place of .i) they modify the entire sentence (bridi), and after the attitudinal beginning bracket « "fu'e" », they continue the emotion for all the statement until the closing bracket « "fu'o" ». Attitudinals can also affect vocatives, the same way a change in vocal tone can greatly affect a simple "hello". For example, « "coi .ui" » means not a simple "Hello.", but something like "Hello! :D". Combined with the attitudinal question-maker "pei", you can easily create more complex « "coi .uipei" » may be thought of as "Hello - are you glad to see me?", while « "co'o .uipei" » may be viewed as "Good-bye - are you happy that we part?". Beyond mere emotion, they can mark prepositional ideas such as desire ("I want this to happen"), obligation ("This needs to happen"), and many other complex functions we take for granted - for example, « "mi citka" » means "I eat", while « ".au mi citka" » means "I want to eat", and « ".e'e mi citka" » means "I am able to eat". Attitudinal indicators and suffixes with examples of modifier use:. Miscellaneous indicators:. See also: Emotions 

Since {ki'e} means approximately "Thank you," it might seem appropriate to respond by saying {fi'i}, seemingly corresponding to the English "You're welcome." However, {fi'i} means "welcome" in the sense of an invitation/acceptance only. In order to respond, {je'e} should be used. "ju'i" can be used as an attention-grabber such as the first word in each of the following utterances (if they started a conversation): "Attention, all shoppers", "So, I just got back from my date with Ryan", "Lo! Listen to the news which I bring". This is the function of the very first word of "Beowulf". It can also return an audience's attention to the speaker/discussion, such as with some usages of "ahem", clearing the throat, or wrapping on a desk. Exercises. In Exercises 1-7, give the most suitable Lojban word(s) to say in each situation. 01. The beginning of a conversation.&lt;br&gt; 02. Begging.&lt;br&gt; 03. During an oath.&lt;br&gt; 04. Inviting a guest into your house.&lt;br&gt; 05. Getting a group of people to listen to your announcement.&lt;br&gt; 06. You desire the floor while someone else has it.&lt;br&gt; 07. Leaving a party. In Exercises 8-43, give the closest Lojban translation of the English sentences. 08. Attention!&lt;br&gt; 09. Listen to me, Victoria.&lt;br&gt; 10. I swear, Benjamin.&lt;br&gt; 11. I don't understand, Christian.&lt;br&gt; 12. I am Pete.&lt;br&gt; 13. Goodbye, Jonathan.&lt;br&gt; 14. Roger, Logan.&lt;br&gt; 15. Wilco.&lt;br&gt; 16. Hello, Rodrigo.&lt;br&gt; 17. Hi, Bertand.&lt;br&gt; 18. Wait! We're not done discussing this!&lt;br&gt; 19. Hark!&lt;br&gt; 20. Agreed.&lt;br&gt; 21. Can I say something, Jose?&lt;br&gt; 22. Just a minute, Jennifer.&lt;br&gt; 23. I'm not done talking, Morgan.&lt;br&gt; 24. Please, Sydney.&lt;br&gt; 25. Over and out.&lt;br&gt; 26. I appreciate that, Chloe.&lt;br&gt; 27. Make yourself at home, Rachel.&lt;br&gt; 28. May I speak, Jasmine?&lt;br&gt; 29. Sorry for interrupting, Sophia, but...&lt;br&gt; 30. I promise, Megan.&lt;br&gt; 31. Thanks, Samuel.&lt;br&gt; 32. Hold on, Nathan.&lt;br&gt; 33. What did you say, Natalie?&lt;br&gt; 34. I interrupt.&lt;br&gt; 35. Greetings, Justin.&lt;br&gt; 36. Uh-huh.&lt;br&gt; 37. Thank you, Dylan.&lt;br&gt; 38. At your service, Austin.&lt;br&gt; 39. I'm listening, Kevin.&lt;br&gt; 40. I understand, Julia.&lt;br&gt; 41. "Houston, we have a problem."&lt;br&gt; 42. Hello! Susan has just left Harold.&lt;br&gt; 43. Hello, Susan! Harold has just left. Answers to exercises. 01. coi&lt;br&gt; 02. pe'u&lt;br&gt; 03. nu'e&lt;br&gt; 04. fi'i&lt;br&gt; 05. ju'i&lt;br&gt; 06. ta'a&lt;br&gt; 07. co'o 08. ju'i&lt;br&gt; 09. ju'i. vektarias.&lt;br&gt; 10. nu'e. bendjamen.&lt;br&gt; 11. je'enai. krestien.&lt;br&gt; 12. mi'e. pit.&lt;br&gt; 13. co'o. janaten.&lt;br&gt; 14. je'e. logen.&lt;br&gt; 15. vi'o&lt;br&gt; 16. coi. radrigos.&lt;br&gt; 17. coi. bertrend.&lt;br&gt; 18. fe'onai&lt;br&gt; 19. ju'i&lt;br&gt; 20. vi'o&lt;br&gt; 21. be'e. xozeis.&lt;br&gt; 22. re'inai. djenefer.&lt;br&gt; 23. mu'onai. morgen.&lt;br&gt; 24. pe'u. sednis.&lt;br&gt; 25. fe'o&lt;br&gt; 26. ki'e. clos.&lt;br&gt; 27. fi'i. reitcel.&lt;br&gt; 28. be'e. djezmen.&lt;br&gt; 29. ta'a. sofias.&lt;br&gt; 30. nu'e. meigen.&lt;br&gt; 31. ki'e. semiul.&lt;br&gt; 32. re'inai. neitan.&lt;br&gt; 33. ke'o. netalis.&lt;br&gt; 34. ta'a&lt;br&gt; 35. coi. jesten.&lt;br&gt; 36. je'e&lt;br&gt; 37. ki'e. delen.&lt;br&gt; 38. fi'i. .asten.&lt;br&gt; 39. re'i. keven.&lt;br&gt; 40. je'e. djulias.&lt;br&gt; 41. doi xustyn. mi'a se nabmi&lt;br&gt; 42. coi do'u la suzyn. la xeirold. puzi cliva&lt;br&gt; 43. coi la suzyn. la xeirold. puzi cliva 

Introduction. Set Theory starts very simply: it examines whether an object "belongs", or does "not belong", to a "set" of objects which has been described in some non-ambiguous way. From this simple beginning, an increasingly complex (and useful!) series of ideas can be developed, which lead to notations and techniques with many varied applications. Definition: Set. The present definition of a set may sound very vague. A set can be defined as an unordered collection of entities that are related because they obey a certain rule. 'Entities' may be anything, "literally": numbers, people, shapes, cities, bits of text, ... etc The key fact about the 'rule' they all obey is that it must be "well-defined". In other words, it must describe clearly what the entities obey. If the entities we're talking about are words, for example, a well-defined rule is:  X is English A rule which is not well-defined (and therefore couldn't be used to define a set) might be:  X is hard to spell Elements. An entity that belongs to a given set is called an element of that set. For example: Set Notation.  formula_1  formula_2  formula_4  formula_5  formula_8  formula_9  formula_10 Note that the use of ellipses may cause ambiguities, the set above may be taken as the set of integers individible by 4, for example. Special Sets. The universal set. The set of all the entities in the current context is called the universal set, or simply the universe. It is denoted by formula_15. The context may be a homework exercise, for example, where the Universal set is limited to the particular entities under its consideration. Also, it may be any arbitrary problem, where we clearly know where it is applied. The empty set. The set containing no elements at all is called the null set, or empty set. It is denoted by a pair of empty braces: formula_16 or by the symbol formula_17. It may seem odd to define a set that contains no elements. Bear in mind, however, that one may be looking for solutions to a problem where it isn't clear at the outset whether or not such solutions even exist. If it turns out that there isn't a solution, then the set of solutions is empty. For example: Operations on the empty set. Operations performed on the empty set (as a set of things to be operated upon) can also be confusing. (Such operations are nullary operations.) For example, the sum of the elements of the empty set is zero, but the product of the elements of the empty set is one (see empty product). This may seem odd, since there are no elements of the empty set, so how could it matter whether they are added or multiplied (since “they” do not exist)? Ultimately, the results of these operations say more about the operation in question than about the empty set. For instance, notice that zero is the identity element for addition, and one is the identity element for multiplication. Special numerical sets. Several sets are used so often, they are given special symbols. Natural numbers. Note that, when we write this set by hand, we can't write in bold type so we write an N in blackboard bold font: formula_22 Integers. In blackboard bold, it looks like this: formula_23 Real numbers. If we expand the set of integers to include all decimal numbers, we form the set of real numbers. The set of reals is sometimes denoted by R. A real number may have a "finite" number of digits after the decimal point (e.g. 3.625), or an "infinite" number of decimal digits. In the case of an infinite number of digits, these digits may: In blackboard bold: formula_24 Rational numbers. Those real numbers whose decimal digits are finite in number, or which recur, are called rational numbers. The set of rationals is sometimes denoted by the letter Q. A rational number can always be written as exact fraction "p"/"q"; where "p" and "q" are integers. If "q" equals 1, the fraction is just the integer "p". Note that "q" may NOT equal zero as the value is then undefined. In blackboard bold: formula_25 Irrational numbers. If a number "can't" be represented exactly by a fraction "p"/"q", it is said to be irrational. Set Theory Exercise 1. Click the link for Set Theory Exercise 1 Relationships between Sets. We’ll now look at various ways in which sets may be related to one another. Equality. Two sets formula_11 and formula_12 are said to be equal if and only if they have exactly the same elements. In this case, we simply write:  formula_28 Note two further facts about equal sets: So, for example, the following sets are all equal:  formula_29 Subsets. If all the elements of a set formula_11 are also elements of a set formula_12, then we say that formula_11 is a subset of formula_12 and we write:  formula_35 For example: In the examples below:  If formula_36 and formula_37, then formula_38  If formula_39 and formula_40, then formula_41  If formula_42 and formula_43, then formula_44 Notice that formula_35 does not imply that formula_12 must necessarily contain extra elements that are not in formula_11; the two sets could be equal – as indeed formula_48 and formula_49 are above. However, if, in addition, formula_12 does contain at least one element that isn’t in formula_11, then we say that formula_11 is a proper subset of formula_12. In such a case we would write:  formula_54 In the examples above:  formula_55 contains ... -4, -2, 0, 2, 4, 6, 8, 10, 12, 14, ... , so formula_56  formula_57 contains $, ;, &amp;, ..., so formula_58 But formula_48 and formula_49 are just different ways of saying the same thing, so formula_61 The use of formula_62 and formula_63; is clearly analogous to the use of &lt; and ≤ when comparing two numbers. Notice also that "every" set is a subset of the "universal set", and the "empty set" is a subset of "every" set. Finally, note that if formula_68 and formula_12 must contain exactly the same elements, and are therefore equal. In other words:  formula_70 Disjoint. Two sets are said to be disjoint if they have no elements in common. For example:  If formula_71 and formula_72, then formula_11 and formula_12 are disjoint sets Venn Diagrams. A "Venn diagram" can be a useful way of illustrating relationships between sets. In a Venn diagram: On the left, the sets "A" and "B" are disjoint, because the loops don't overlap. On the right "A" is a subset of "B", because the loop representing set "A" is entirely enclosed by loop "B". &lt;br clear="all"&gt; Venn diagrams: Worked Examples. "Example 1" "Fig. 3" represents a Venn diagram showing two sets "A" and "B", in the general case where nothing is known about any relationships between the sets. Note that the rectangle representing the universal set is divided into four regions, labelled "i", "ii", "iii" and "iv". What can be said about the sets "A" and "B" if it turns out that: &lt;br clear="all"&gt; (a) If region "ii" is empty, then "A" contains no elements that are not in "B". So "A" is a subset of "B", and the diagram should be re-drawn like "Fig 2" above. (b) If region "iii" is empty, then "A" and "B" have no elements in common and are therefore disjoint. The diagram should then be re-drawn like "Fig 1" above. "Example 2" (a) The diagram in "Fig. 4" shows the general case of three sets where nothing is known about any possible relationships between them. (b) The rectangle representing U is now divided into 8 regions, indicated by the Roman numerals "i" to "viii". (c) Various combinations of empty regions are possible. In each case, the Venn diagram can be re-drawn so that empty regions are no longer included. For example: "Example 3" The following sets are defined: Using the two-stage technique described below, draw a Venn diagram to represent these sets, marking all the elements in the appropriate regions. The technique is as follows: Don't begin by entering the elements of set "A", then set "B", then "C" – you'll risk missing elements out or including them twice! "Solution" After drawing the three empty loops in a diagram looking like "Fig. 4" (but without the Roman numerals!), go through each of the ten elements in U - the numbers 1 to 10 - asking each one three questions; like this: First element: 1 A 'no' to all three questions means that the number 1 is outside all three loops. So write it in the appropriate region (region number "i" in "Fig. 4"). Second element: 2 Yes, yes, no: so the number 2 is inside "A" and "B" but outside "C". Goes in region "iii" then. ...and so on, with elements 3 to 10. The resulting diagram looks like "Fig. 5". &lt;br clear="all"&gt; The final stage is to examine the diagram for empty regions - in this case the regions we called "iv", "vi" and "vii" in "Fig. 4" - and then re-draw the diagram to eliminate these regions. When we've done so, we shall clearly see the relationships between the three sets. So we need to: The finished result is shown in "Fig. 6". The regions in a Venn Diagram and Truth Tables. Perhaps you've realized that adding an additional set to a Venn diagram "doubles" the number of regions into which the rectangle representing the universal set is divided. This gives us a very simple pattern, as follows: It's not hard to see why this should be so. Each new loop we add to the diagram divides each existing region into two, thus doubling the number of regions altogether. But there's another way of looking at this, and it's this. In the solution to "Example 3" above, we asked three questions of each element: "Are you in A? Are you in B?" and "Are you in C?" Now there are obviously two possible answers to each of these questions: "yes" and "no". When we "combine" the answers to three questions like this, one after the other, there are then 23 = 8 possible sets of answers altogether. Each of these eight possible combinations of answers corresponds to a different region on the Venn diagram. The complete set of answers resembles very closely a "Truth Table" - an important concept in "Logic", which deals with statements which may be "true" or "false". The table on the right shows the eight possible combinations of answers for 3 sets "A", "B" and "C". You'll find it helpful to study the patterns of Y's and N's in each column. Set Theory Exercise 2. Click link for Set Theory Exercise 2 Operations on Sets. Just as we can combine two numbers to form a third number, with operations like 'add', 'subtract', 'multiply' and 'divide', so we can combine two sets to form a third set in various ways. We'll begin by looking again at the Venn diagram which shows two sets "A" and "B" in a general position, where we don't have any information about how they may be related. &lt;br clear="all"&gt; The first two columns in the table on the right show the four sets of possible answers to the questions "Are you in A?" and "Are you in B?" for two sets "A" and "B"; the Roman numerals in the third column show the corresponding region in the Venn diagram in "Fig. 7". Intersection. Region "iii", where the two loops overlap (the region corresponding to 'Y' followed by 'Y'), is called the "intersection" of the sets "A" and "B". It is denoted by "A" ∩ "B". So we can define intersection as follows: For example, if "A" = {1, 2, 3, 4} and "B" = {2, 4, 6, 8}, then "A" ∩ "B" = {2, 4}. We can say, then, that we have combined two sets to form a third set using the "operation of intersection". Union. In a similar way we can define the "union" of two sets as follows: The union, then, is represented by regions "ii", "iii" and "iv" in "Fig. 7". You'll see, then, that in order to get into the intersection, an element must answer 'Yes' to "both" questions, whereas to get into the union, "either" answer may be 'Yes'. The ∪ symbol looks like the first letter of 'Union' and like a cup that will hold a lot of items. The ∩ symbol looks like a spilled cup that won't hold a lot of items, or possibly the letter 'n', for i'n'tersection. Take care not to confuse the two. Difference. This is written "A" - "B", or sometimes "A \ "B". The elements in the difference, then, are the ones that answer 'Yes' to the first question "Are you in A?", but 'No' to the second "Are you in B?". This combination of answers is on row 2 of the above table, and corresponds to region "ii" in "Fig.7". Complement. So far, we have considered operations in which "two" sets combine to form a third: "binary" operations. Now we look at a "unary" operation - one that involves just "one" set. Clearly, this is the set of elements that answer 'No' to the question "Are you in A?". Notice the spelling of the word "complement": its literal meaning is 'a complementary item or items'; in other words, 'that which completes'. So if we already have the elements of "A", the complement of "A" is the set that "completes" the universal set. Summary. &lt;br clear="all"&gt; Cardinality. Finally, in this section on "Set Operations" we look at an operation on a set that yields not another set, but an integer. Generalized set operations. If we want to denote the intersection or union of "n" sets, "A"1, "A"2, ..., "A""n" (where we may not know the value of "n") then the following "generalized set notation" may be useful: In the symbol formula_90 "A""i", then, "i" is a variable that takes values from 1 to "n", to indicate the repeated intersection of all the sets "A"1 to "A""n". Set Theory Exercise 3. Click link for Set Theory Exercise 3 



Intro: What Is Biochemistry? Biochemistry is the study of the chemistry of, and relating to, biological organisms. It forms a bridge between biology and chemistry by studying how complex chemical reactions and chemical structures give rise to life and life's processes. Biochemistry is sometimes viewed as a hybrid branch of organic chemistry which specializes in the chemical processes and chemical transformations that take place inside of living organisms, but the truth is that the study of biochemistry should generally be considered neither fully "biology" nor fully "chemistry" in nature. Biochemistry incorporates everything in size between a molecule and a cell and all the interactions between them. The aim of biochemists is to describe in molecular terms the structures, mechanisms and chemical processes shared by all organisms, providing organizing principles that underlie life in all its diverse forms. Biochemistry essentially remains the study of the structure and function of cellular components (such as enzymes and cellular organelles) and the processes carried out both on and by organic macromolecules - especially proteins, but also carbohydrates, lipids, nucleic acids, and other biomolecules. All life forms alive today are generally believed to have descended from a single proto-biotic ancestor, which could explain why all known living things naturally have similar biochemistries. Even when it comes to matters which could appear to be arbitrary - such as the genetic code and meanings of codons, or the "handedness" of various biomolecules - it is irrefutable fact that all marine and terrestrial living things demonstrate certain unchanging patterns throughout every level of organization, from family and phylum to kingdom and clade. Biochemistry is, most simply put, the chemistry of life. 

Catalysis refers to the acceleration of the rate of a chemical reaction by a substance, called a catalyst, that is itself unchanged by the overall reaction. Catalysis is crucial for any known form of life, as it makes chemical reactions happen much faster than they would "by themselves", sometimes by a factor of several million times. A common misunderstanding is that catalysis "makes the reaction happen", that the reaction would not otherwise proceed without the presence of the catalyst. However, a catalyst cannot make a thermodynamically unfavorable reaction proceed. Rather, it can only speed up a reaction that is already thermodynamically favorable. Such a reaction in the absence of a catalyst would proceed, even without the catalyst, although perhaps too slowly to be observed or to be useful in a given context. Catalysts accelerate the chemical reaction by providing a lower energy pathway between the reactants and the products. This usually involves the formation of an intermediate, which cannot be formed without the catalyst. The formation of this intermediate and subsequent reaction generally has a much lower activation energy barrier than is required for the direct reaction of reactants to products. Catalysis is a very important process from an industrial point of view since the production of most industrially important chemicals involve catalysis. Research into catalysis is a major field in applied science, and involves many fields of chemistry and physics. Two types of catalysis are generally distinguished. In homogeneous catalysis the reactants and catalyst are in the same phase. For example acids (H+ ion donors) are common catalysts in many aqueous reactions. In this case both the reactants and the catalysts are in the aqueous phase. In heterogeneous catalysis the catalyst is in a different phase than the reactants and products. Usually, the catalyst is a solid and the reactants and products are gases or liquids. In order for the reaction to occur one or more of the reactants must diffuse to the catalyst surface and adsorb onto it. After reaction, the products must desorb from the surface and diffuse away from the solid surface. Frequently, this transport of reactants and products from one phase to another plays a dominant role in limiting the rate of reaction. Understanding these transport phenomena is an important area of heterogeneous catalyst research. Enzymes. Enzyme (from Greek, "in ferment") are special protein molecules whose function is to facilitate or otherwise accelerate most chemical reactions in cells. They are simply biological catalysts. Most enzymes are proteins, although a few are catalytic RNA molecules called ribozymes. Many chemical reactions occur within biological cells, but without catalysts most of them happen too slowly in the test tube to be biologically relevant. Enzymes can also serve to couple two or more reactions together, so that a thermodynamically favorable reaction can be used to "drive" a thermodynamically unfavorable one. One of the most common examples is enzymes which use the dephosphorylation of ATP to drive some otherwise unrelated chemical reaction. Chemical reactions need a certain amount of activation energy to take place. Enzymes can increase the reaction speed by favoring or enabling a different reaction path with a lower activation energy (Fig. 1), making it easier for the reaction to occur. Enzymes are large globular proteins that catalyze (accelerate) chemical reactions. They are essential for the function of cells. Enzymes are very specific as to the reactions they catalyze and the chemicals (substrates) that are involved in the reactions. Substrates fit their enzymes like a key fits its lock (Fig. 2). Many enzymes are composed of several proteins that act together as a unit. Most parts of an enzyme have regulatory or structural purposes. The catalyzed reaction takes place in only a small part of the enzyme called the active site, which is made up of approximately 2 - 20 amino acids. The substrates (A and B) need a large amount of energy ("E"1) to reach the intermediate state A...B, which then reacts to form the end product (AB). The enzyme (E) creates a microenvironment in which A and B can reach the intermediate state (A...E...B) more easily, reducing the amount of energy needed ("E"2). As a result, the reaction is more likely to take place, thus improving the reaction speed. Enzymes can perform up to several million catalytic reactions per second. To determine the maximum speed of an enzymatic reaction, the substrate concentration is increased until a constant rate of product formation is achieved (Fig. 3). This is the maximum velocity ("V"max) of the enzyme. In this state, all enzyme active sites are saturated with substrate. This was proposed in 1913 by Leonor Michaelis and Maud Menten. Since the substrate concentration at Vmax cannot be measured exactly, enzymes are characterized by the substrate concentration at which the rate of reaction is half its maximum. This substrate concentration is called the Michaelis-Menten constant ("K"M). Many enzymes obey Michaelis-Menten kinetics. The speed "V" means the number of reactions per second that are catalyzed by an enzyme. With increasing substrate concentration [S], the enzyme is asymptotically approaching its maximum speed "V"max, but never actually reaching it. Because of that, no [S] for "V"max can be given. Instead, the characteristic value for the enzyme is defined by the substrate concentration at its half-maximum speed ("V"max"/2"). This KM value is also called Michaelis-Menten constant. Several factors can influence the reaction speed, catalytic activity, and specificity of an enzyme. Besides de novo synthesis (the production of more enzyme molecules to increase catalysis rates), properties such as pH or temperature can denature an enzyme (alter its shape) so that it can no longer function. More specific regulation is possible by posttranslational modification (e.g., phosphorylation) of the enzyme or by adding cofactors like metal ions or organic molecules (e.g., NAD+, FAD, CoA, or vitamins) that interact with the enzyme. Allosteric enzymes are composed of several subunits (proteins) that interact with each other and thus influence each other's catalytic activity. Enzymes can also be regulated by competitive inhibitors (Fig. 4) and non-competitive inhibitors and activators (Fig. 5). Inhibitors and activators are often used as medicines, but they can also be poisonous. A competitive inhibitor fits the enzyme as well as its real substrate, sometimes even better. The inhibitor takes the place of the substrate in the active center, but cannot undergo the catalytic reaction, thus inhibiting the enzyme from binding with a substrate molecule. Some inhibitors form covalent bonds with the enzyme, deactivating it permanently (suicide inhibitors). In terms of the kinetics of a competitive inhibitor, it will increase Km but leave Vmax unchanged. Non-competitive inhibitors/activators (I) do not bind to the active center, but to other parts of the enzyme (E) that can be far away from the substrate (S) binding site. By changing the conformation (the three-dimensional structure) of the enzyme (E), they disable or enable the ability of the enzyme (E) to bind its substrate (S) and catalyze the desired reaction. The noncompetitive inhibitor will lower Vmax but leave Km unchanged. An "uncompetitive" inhibitor will only bind to the enzyme-substrate complex forming an enzyme-substrate-inhibitor (ESI) complex and cannot be overcome by additional substrate. Since the ESI is nonreactive, Vmax is effectively lowered. The uncompetitive inhibitor will in turn lower the Km due to a lower concentration of substrate needed to achieve half the maximum concentration of ES. Several enzymes can work together in a specific order, creating metabolic pathways (e.g., the citric acid cycle, a series of enzymatic reactions in the cells of aerobic organisms, important in cellular respiration). In a metabolic pathway, one enzyme takes the product of another enzyme as a substrate. After the catalytic reaction, the product is then passed on to another enzyme. The end product(s) of such a pathway are often non-competitive inhibitors (Fig. 5) for one of the first enzymes of the pathway (usually the first irreversible step, called "committed step"), thus regulating the amount of end product made by the pathway (Fig. 6). Enzymes are essential to living organisms, and a malfunction of even a single enzyme out of approximately 2,000 present in our bodies can lead to severe or lethal illness. An example of a disease caused by an enzyme malfunction in humans is phenylketonuria (PKU). The enzyme phenylalanine hydroxylase, which usually converts the essential amino acid phenylalanine into tyrosine does not work, resulting in a buildup of phenylalanine that leads to mental retardation. Enzymes in the human body can also be influenced by inhibitors in good or bad ways. Aspirin, for example, inhibits an enzyme that produces prostaglandins (inflammation messengers), thus suppressing pain. But not all enzymes are in living things. Enzymes are also used in everyday products such as biological washing detergents where they speed up chemical reactions, (to get your clothes clean). Digestive and Metabolic Enzymes. In the previous section we have been talking about the digestive enzymes, both the ones produced by the body, such as salivary amylase, and the food enzymes. Their primary role is for the digestion of food. Another class of enzymes is called metabolic enzymes. Their role is to catalyze chemical reactions involving every process in the body, including the absorption of oxygen. Our cells would literally starve for oxygen even with an abundance of oxygen without the action of the enzyme, cytochrome oxidase. Enzymes are also necessary for muscle contraction and relaxation. The fact is, without both of these classes of enzymes, (digestive and metabolic,) life could not exist. Digestive enzymes function as biological catalysts in which it helps to breakdown carbohydrates, proteins, and fats. On the other hand, metabolic enzymes function as a remodel of cells. Digestion of food has a high priority and demand for enzymes; digestive enzymes get priority over metabolic enzymes. Any deficiency in metabolic enzyme can lead to over work, which could lead to enlarge organs in order to perform the increased workload. The result is unhealthy and could cause enlarged heart or pancreas. The deficiencies of metabolic enzymes can have a tremendous impact on health. As we grow older enzyme level decline and the efficiency in the body decline. Enzyme naming conventions. By common convention, an enzyme's name consists of a description of what it does, with the word ending in "-ase". Examples are alcohol dehydrogenase and DNA polymerase. Kinases are enzymes that transfer phosphate groups. The International Union of Biochemistry and Molecular Biology has developed a nomenclature for enzymes, the EC numbers; each enzyme is described by a sequence of four numbers, preceded by "EC". The first number broadly classifies the enzyme based on its mechanism: Some other important enzymes are: Protease: breaks the protein into amino acids in high acidity environments such as stomach, pancreatic and intestinal juices. Act on bacteria, viruses and some cancerous cells. Amylase: Break complex carbohydrates such as starch into simpler sugars (dextrin and maltose). It found in the intestines, pancreas and also in salivary glands. Lipase: breaks down fats and some fat soluble vitamins (A,E,K, and D). helpful in treating cardiovascular diseases. Cellulase: break down cellulose that found in fruits, grains, and vegetables. It increases the nutritional values of vegetables, and fruits. Pectinase: break down pectin that found in citrus fruits, carrots, beets, tomatoes, and apple. Antioxidants: protect from free radical negative effect that can damage cell in the body. Cathepsin: break animal protein down. Lactase: break down lactose that found in milk products. the production of lactase decrease with age. Invertase: assimilate sucrose that can contribute to digestive stress if not digested properly. Papain: break down protein and help the body in digestion. Bromelain: Break proteins that found in plants and animals. it could help the body to fight cancer and treat inflammation. Glucoamylase: break down maltose that found in all grains in to two glucose molecules. 

This is a guide to HTML, a standard for web pages. A text editor and a web browser is all you need to create web pages, view your handiwork, and share information with others all over the world. This book covers simple HTML syntax. For dynamic behavior in websites, see the JavaScript wikibook. Another separate book covers Cascading Style Sheets (CSS) which handle overall look and styling, but the present book addresses CSS briefly. Additionally, XHTML has its own textbook. Other Wikibooks.  

http://www.cde.ca.gov/re/pn/fd/documents/hist-social-sci-frame.pdf http://www.cde.ca.gov/be/st/ss/hstgrades9through12.asp 

To study Spanish or any non-native language you should determine your goals. Many people study a language in school to fulfill a requirement. Others do it generally as a "broadening" experience. Still others do it so that they can speak with others in their own language. Number one tip: study every day. To speak any language fluently and be able to survive in the language among native speakers, it is probably necessary to immerse yourself in the language for an extended period such as 9 months to a year. Short of that it is best to study intensively every day for an extended period. Studying every day is much more effective than studying every other day. Studying only once a week is a sure way to forget everything by the end of the week and be eternally relearning the same little bit of vocabulary. The problem is that your mind will forget much of what you learn in the course of one day and by the end of two days it is even more. If you study every day, even for a little while, it will mean that you can focus more of your time learning new material and waste less time relearning old material. This is one of the secrets to learning a language: study it every day. Tip for studying new words: Have your own vocabulary. It can be some note, every page is e.g. 10 words. And every day study 10 new words, 10 words from the yesterday, 10 words from the day before yesterday, 10 words from the week before, 10 from 2 week before, 10 from 3 week before and 10 from 4 week before. It's overall 70 words. It want about half or 1 hour studying. 10 new words want more time than then repeat 60 others and you have confidence that you will know for a long time. Spaced learning software such as anki does a similar thing by scheduling reviews depending on how well you do. Using web resources. Learning via the web offers a lot of benefits that a standard textbook can not provide. One example are audio files which we already include in some of our lessons. Another benefit are online dictionaries that allow you to quickly look up vocabulary. Check out our Web Resources page to find out more. Using media resources. Once you have learned the basics of Spanish, try reading or viewing material in spanish. A newspaper, a magazine or a children's book is a great place to start. Online retailers, such as Amazon.com, have international storefronts that sell books in foreign languages. Watching Spanish television or Spanish movies will help you practice listening. Playing your favorite DVD or television show with Spanish close captioning on or with the Spanish soundtrack will help you learn new vocabulary and will assist in putting items into context. Areas where a significant number of people speak Spanish should have reading material available at local supermarkets and bookstores, and may have some local television or radio stations in Spanish. If you cannot find Spanish media where you live, the Internet has online editions of many Spanish newspapers and magazines available. Keep a dictionary nearby if you have trouble with unfamiliar vocabulary. The dictionary can be a very effective study tool. Each time you look up a new word, place a dot or symbol near the word. If you forget a word and need to look it up again, add another dot near the word. Pretty soon, you will have a collection of words with two, three and four dots. This marking system helps you figure out which words are either important to remember (because you use them frequently) or which ones are difficult to remember. During your vocabulary review, you can flip through your dictionary and find important words to practice. Regular use of a dictionary has the added benefit of putting things into context. Many dictionaries have tidbits and trivia added to the definitions, to help you see how components of language fit together. In addition, every time you open your dictionary, you will see other Spanish words, while searching for your target word and this has the effect of reminding you of other words. Use it. The best way to learn Spanish is by using it with other speakers. Do not be discouraged by your bad accent and choppy flow; these are things that come with being a beginner. People will understand that you are learning and will appreciate you applying the skills that you have acquired. So go ahead—say something. Ask your neighbor "¿Dónde está el baño?" Look around your local community. Often times coffee shops, libraries, or other local organization sponsor a "Spanish club" for adults. Both native and non-native speakers attend. This is a great way to use Spanish by speaking with others (especially native speakers). Another great way to practice your new language skills is to travel to the country where the language is spoken. Most major cities in Spanish speaking countries host a number of Spanish language schools for foreigners. These schools offer language classes at all levels by hour, week, or month of classes. Also, they normally will arrange for their students to live with a local host family, which provides another source for practicing the Spanish language with native speakers. 

Grammar - Questions. Unlike English, yes/no questions in Spanish are not usually formed by switching the position of subject and verb (if the subject is explicit). To recognize a sentence as affirmative or as a question one must pay attention to the intonation pattern. Unlike English, Spanish uses a reversed question mark (¿) at the beginning of a question: become For other type of questions Spanish uses the following question words (note that all of them have an accent in the word): Here are some Spanish sentences where specific question words are used: Questions can also be posed within a sentence: Exercise: Questions Grammar - Possessive Adjectives. Like English, the Spanish possessive adjectives differ depending on the person they are referring to. Unlike English, the possessive article also changes depending on the number of items that one possesses (for example: mi libro = "my book", mis libros = "my books"). It can also change depending on the gender of the item (for example: nuestro perro = "our dog", nuestra casa = "our house"). The following table summarizes all Spanish possessive adjectives: Exercise: Possessive Adjectives Grammar - Comparisons. Equality. Spanish uses three slightly different constructions for comparisons of equality. One for comparing verbs, one for comparing nouns and one for comparing adjectives/adverbs. The following examples show the three different possibilities: When comparing nouns, the ending of tanto will be modified to tanta, tantos, or tantas in order to match gender and quantity of the noun. The general pattern for comparisons of equality is the following: Inequality. For comparisons of inequality, Spanish uses the same form for both nouns and adjectives/adverbs. There are two types of inequalities: más ... que ("more than") and menos ... que ("less than"): In general: Superlatives. Superlatives in Spanish are similar to comparisons of inequality: They use más for "the most", menos for "the least". Then follows the adjective and finally there is a preposition (de): Note that in some cases (la más inteligente) you can just write the article and omit the noun. The general pattern for Spanish superlatives is: Exercise: Comparisons 

=Chordates= The phylum Chordata includes three subphyla. These include vertebrates and invertebrate chordates. Characteristics. Notochord: the rod-shaped supporting axis found in the dorsal part of the embryos of all chordates, including vertebrates Flexible, non-collapsible rod dorsal to the gut/coelom and below the nervous system, hydrostatic, fluid wrapped in tough connective tissue. As bone does not compact, muscles tensed on one side result in movement instead of shortening the animal. This allows much better locomotion than do cilia for larger animals in water, a crucial victory for later success. Pharyngeal slits: Slits in the pharynx originally used to gather food, water enters the mouth, passes through pharynx and out gill-like slits, passing through a cavity called an antrium and then outside. In humans, present only in embryo. Dorsal nerve cord: A neural tube dorsal to the notochord Postanal tail: Elongation of the body and notochord, nerve cord and muscles past anus into tail, early locomotive function led to success. Non-synapomorphic characteristics (not limited to chordates): Subphylum Urochordata. The tunicates are located in this subphylum. Along with the subphylum Cephalochordata, these two subphyla make up the invertebrate chordates. Only the tunicate larvae have notochords, nerve cords, and postanal tails. Most adult tunicates are sessile, filter-feeders which retain their pharyngeal slits. Adult tunicates also develop a sac, called a tunic, which gives tunicates their name. Cilia beating within the turnicate cause water to enter the incurrent siphon. The water enters the body, passes through the pharyngeal slits, and leaves the body through the excurrent siphon. Undigested food is removed through the anus. Tunicates are hemaphrodites and can reproduce asexually through budding.In urochordates notochord is confined to larval tail.These lack cranium. These have an open type of circulatory system.Excretion is by neural gland, nephrocytes.there are two siphons through which water enters and exit.they have a tubular heart.they have a tough outer covering.example acidia Subphylum Cephalochordata. The lancelets are located in this subphylum. Along with the subphylum Urochordata, these two subphyla make up the invertebrate chordates. Lancelets receive their name from their bladelike shape. They resemble fish but they are actually scaleless chordates only a few centimeters long. They spend most of their time buried in the sand with their mouths protruding. Fossils of lancelets have been found to be over 550 million years old. Dropped out sessile stage, what was the larval stage is now sexually reproductive. Includes Branchiostoma (“amphioxus”). Subphylum Vertebrata. (Vertebra from Latin vertere, to turn). Characterized by separate bones or cartilage blocks firmly joined as a backbone. The backbone supports and protects a dorsal nerve cord. Vertebrates have tissues which are organized into organs which in turn are organized into organ systems. All vertebrates share the following characteristics: - segmentation - a true coelom - bilateral symmetry - cephalization - a backbone - a bony skull - a closed circulatory system - chambered heart - two pairs of jointed appendages - tissues organized into organs Vertebrate Organ Systems: - Nervous System - Circulatory System - Digestive System - Respiratory System - Reproductive System - Excretory System What evolutionary relationship could we imagine between sessile echinoderms and the higher chordate animals? Paedomorphic (child-form) hypothesis: basically, evolution of sexual reproduction in what had previously been a larval life stage, or the retention of at least one juvenile characteristic into the adult (adult = sexually reproducing) stage. Some scientists believe that this occurred in a proto-chordate animal lineage. Maybe chordates (and vertebrates) arose from sessile (attached) ancestors. Selection in these proto-chordates maybe began to favor more time in the larval stage, as feeding was more successful or mortality lower in this stage. As larvae got bigger physics shows that the cilia become less efficient for locomotion, favoring the undulating motion allowed by a notochord. Is this hypothesis crazy? A similar example of this today is Epemeroptera, the mayfly, which has almost abandoned its adult stage. Its one-year lifespan is mostly larval with just a brief day of reproduce-and-die as an adult, which doesn’t even have usable mouthparts. Tunicate (sea squirt) larva has all four chordate characteristics, although adult sessile (“attached”). Class Agnatha. "jawless fish" Appeared approximately 500 million years ago and dominated the oceans for about 100 million years. The first group of fish to appear. They had neither jaws, paired fins, nor scales, but they were the first organisms with a backbone. Class Acanthodia. "spiny fish" Appeared about 430 million years ago. An extinct class of fish that developed jaws with bony edges. They had internal skeletons made of cartilage and some bone. Class Placodermi. Appeared about 410 million years ago, dominated the sea for about 50 million years. An extinct class of fisive heads. Class Chondrichthyes. "cartilaginous fish" Appeared about 400 million years ago with bony fish. Includes sharks, skates and rays, and chimaeras. Their skeletons are made of cartilage strengthened by the mineral calcium carbonate. The main characteristics and distinguishing features of this class: - gills - single-loop blood circulation - vertebral column - presence of placoid scales on their bodies - internal skeleton of cartilage - paired, fleshy pectoral and pelvic fins - asymmetrical tail fin prevents sinking - fatty liver provides neutral buoyancy - visceral clefts present as separate and distinct gills - no external ear - oviparous - internal fertilization - ectoderms - cold blooded Class Osteichthyes. "bony fish" Appeared about 400 million years ago with cartilaginous fish. Includes about 95% of today's fish species. Subclass Sarcopterygii. fleshy-finned fishes. Fins have bones and muscles, homologous to our limbs. Order Dipnoi. lung fishes, two groups isolated when continents separated Order Crossopterygii. includes coelacanths and rhipodistians, gave rise to amphibians, had lungs which evolved into a swim bladder in bony fishes, and labyrinthodont teeth, characterized by complex folding of enamel. • Skeleton made of bone, jaws, fins, most with scales, two-chambered heart. Class Amphibia. means “both lives”, aquatic larvae, terrestrial adult Amphibians: - Legs - Lungs - Double-Loop Circulation - Partially Divided Heart - Cutaneous Respiration (Breathes through Skin) Order Salientia. frogs (jumping) (aka Anura) Order Urodela. salamanders (tailed) Labyrinthodont amphibians: oldest known amphibians, inherited characteristic teeth from crossopterygii ancestor, had stocky, aquatic larvae. Amphibians have limbs instead of fins. Girdles and vertebral column now more substantial and connected, support body on legs. Lisamphybia: no scales, “smooth”, eggs with no shell, laid in water (water-reliant). Amphibians gave rise to cotylosaurs, from which arose dinosaurs, turtles, lizards, and therapsids. Class Reptilia. amniotic egg allowed freedom from water, shelled egg. (Amnion: protection). Reptiles have four extra-embryonic membranes: Reptiles are cold-blooded, or ectothermic, meaning that their heat come from their environment. Sometimes defined as all amniotes that are not birds or mammals. Reptiles can be classified by skull structure into four groups: Refers to number of holes in the skull. Cotylosaurs had Anapsid skull Dermatocranium: from bony outer skull structure, precursor to human cranium. Subclass Testudinata. turtles, terrapins Subclass Diapsida. dinosaurs, snakes, most stuff Subclass Diapsida. includes Ichthyosaurs, marine reptiles convergent on dolphins; Plesiosaurs, ancient sea monsters; Squamates, including lizards and snakes; and Thecodonts, which gave rise to Dinosaurs: broken into two groups, based on hip structure Crocodilians: come from archosaurs, the only extant (still living today) archosaur descendant. Ancestrally bipedal, secondarily quadripedal. Synapsids: refers to joined (Greek syn-, together with) parts of skull. Led eventually to mammals. Synapsid pelycosaur » therapsid » mammals Pelycosaur: Sail-backed dinosaur, legs not spread out like lizard but more pillar-like and under body, allowing greater activity and competence in motion, pendulum like rather than constant push-up. Teeth differentiated into different types, for pre-processing of food needed by higher metabolism. Skull changes, bone histology, suggestions of warm-bloodedness. Class Aves. arose late Jurassic, early Cretaceous. Feathers, skeleton modified for flight. Feathers: epidermal derivative, made of keratin (like fingernails). Carpometacarpis: bears primary flight feathers, parallel to hand parts. Keeled sternum: breastbone, powerful one needed to support flight muscles. Strong, light, occasionally hollow bones. All birds lay eggs (as contrasted to reptiles, which have developed live birthing over 100 independent times.) Why are there no live-bearing birds? Early birds had teeth, lost them. With mammals, only exothermic animals. Archaeopteryx: “ancient wing”, Jurassic bird-reptile, very dinosaur-like. Good fossils found in Zolenhoffen, German sandstone mine with fine sand, shows feathers clearly, found shortly after Darwin’s publication and used to support his hypothesis. Thick, heavy bones and no sternum, bony tail, not a good flyer but did have primary flight feathers. Archaeornithes: includes archaeopteryx. Paleognathae: gave rise to Australian flightless birds. Neognathae: remaining live birds. Class Mammalia. Two unique characteristics, or synapomorphies: Three skeletal characteristics (fossilize) Mammals basically have a synapsid skull design inherited from ancestor Non diagnostic characteristics (not unique to mammals): Subclass Protheria. monotremes (Greek mon-, one; and trema, hole), or egg-laying mammals, have one opening for excretion and urination. Subclass Theria. Metatheria: Marsupials (opossum, kangaroo…) Eutheria: Placental mammals (all common mammals) Marsupium: (from Greek marsypion, purse or pouch). Gestation period much shorter than in Eutherian mammals, but after leaving the uterus the tiny offspring crawls into a pouch where it completes development latched onto a teat. Recent molecular (read: genetic) evidence suggests that two different mammal groups may have developed live-bearing ability separately. Instead of being a “rough draft” for placental-style live bearing, perhaps the marsupial pouch approach is another solution to the same problem. Advantage: in tough times the parent can pitch out the offspring and increase its own chance of survival. 

General Biology =Vertebrate tissues = Definition of tissue: an aggregation of cells, usually of the same kind, organized to perform a common function. A group of similar cells organized into a structural and functional unit (Raven/Johnson). Primary categories of adult tissues (non-dogmatic categories of convenience): 

&lt;br&gt; &lt;br&gt; 

Introduction to animal phyla. There currently are almost 40 recognized phyla. Phylum — Number of Species — Common Name Phylum Porifera. The name Porifera means "pore-bearing". This phylum is commonly called sponges. The number of species is estimated to be between 5,000 and 10,000. All are aquatic and almost all are marine. Animals in this phyla have no true tissues, which means, for example, that they have no nervous system or sense organs. Although sponges are multicellular, they are described as being essentially at a cellular level of organization. They are sessile as adults, but have a free swimming larva. Their bodies are porous. They are filter feeders; water flows in through many small openings (ostia), and out through fewer, large openings (osculum). They have inner and outer cell layers, and a variable middle layer. The middle layer often is gelatinous with spiny skeletal elements (called spicules) of silica or calcium carbonate, and fibres made of spongin (a form of collagen). Choanocytes are flagellated cells lining the inside of the body that generate a current, and trap and phagocytize food particles. Their cells remain totipotent, or developmentally flexible: they can become any type of cell at any point in the sponge's development. This allows for the great regenerative power sponges have. Sponges are an ancient group, with fossils from the early Cambrian (ca. 540 mya) and possibly from the Precambrian. Sponges often are abundant in reef ecosystems. They somehow are protected from predators (spicules? bad taste?). Many organisms are commensals of sponges, living inside them. Some sponges harbor endosymbiotic cyanobacteria or algae (dinoflagellates, a.k.a. "zooxanthellae"). Phylum Cnidaria. See text pages 886 - 889. Name comes from the Greek knide- meaning "nettle". This phylum formerly known as phylum coelenterata consists of the jellyfish, hydra, sea anemones, corals, sea pens, sea wasps, sea whips and box jellyfish. There are about 9,000 species. Almost all are marine. This is another ancient group, with fossils perhaps reaching back to 700 mya. Cnidarians exhibit radial symmetry. Their basic body plan is a sac with a gastrovascular cavity, or a central digestive system. They have one opening, which serves as both mouth and anus. The body wall has an outer ectoderm, an inner endoderm, and a variable undifferentiated middle layer called mesoglea or mesenchyme that may be jelly-like. The mesoglea is NOT considered to be true mesoderm and so the Cnidaria are described as diploblastic. Tentacles usually extend from the body wall around the mouth/anus. There are two basic body plans: the polyp and the medusa. The polyp is sessile and attaches to substrate by the aboral end (i.e., the end away from the mouth). The medusa ("jellyfish") is a floating form, and looks like an upside-down version of the polyp. Some cnidarians only have the polyp stage, some have only the medusa stage, and others have both. The typical life cycle of a cnidarian involves what is called "alternation of generations": an alternation between an asexual polyp stage and a sexual medusa stage. The tentacles are armed with cnidae (or nematocysts), small intracellular "harpoons" that function in defense and prey capture. When fired, the cnidae deliver a powerful toxin that in some cases is dangerous to humans. The phylum is named after the cnidae. Cnidarians have no head, no centralized nervous system, and no specialized organs for gas exchange, excretion, or circulation. They do have a "nerve net" and nerve rings (in jellyfish). Many cnidarians have intracellular algae living within them in a mutualistic symbiotic relationship (Dinoflagellates = zooxanthellae). This combination is responsible for much of the primary productivity of coral reefs. There are three main classes in the phylum Phylum Platyhelminthes. See text pages 890 - 893. Name means "flat worm" Most members of this phylum are parasitic (flukes and tapeworms), but some are free living (e.g., planaria). There are about 20,000 species. They are dorsoventrally compressed (i.e., "flat"). Animals in this phylum are acoelomate, triploblastic, bilaterally symmetrical, and unsegmented. Platyhelminths have a simple anterior "brain" and a simple ladder-like nervous system. Their gut has only one opening. Flatworms have NO circulatory or gas exchange systems. They do have simple excretory/osmoregulatory structures (protonephridia or "flame cells"). Platyhelminths are hermaphroditic, and the parasitic species often have complex reproductive (life) cycles. There are four main classes of platyhelminths: Phylum Rotifera. See text page 900 The Rotifers. The name means "wheel bearing," a reference to the corona, a feeding structure (see below). They are triploblastic, bilaterally symmetrical, and unsegmented. They are considered pseudocoelomates. Most less than 2 mm, some as large as 2 – 3 mm. Rotifers have a three part body: head, trunk, and foot. The head has a ciliary organ called the corona that, when beating, looks like wheels turning, hence the name of the phylum. The corona is a feeding structure that surrounds the animal's jaws. The gut is complete (i.e., mouth &amp; anus), and regionally specialized. They have protonephridia but no specialized circulatory or gas-exchange structures. Most live in fresh water, a very few are marine or live in damp terrestrial habitats. They typically are very abundant. There are about 2,000 species. Parthenogenesis, where females produce more females from unfertilized but diploid eggs, is common. Males may be absent (as in bdelloid rotifers) or reduced. When males are present, sexual and asexual life cycles alternate. Males develop from unfertilized haploid eggs and are haploid. Males produce sperm by mitosis which can fertilize haploid eggs, yielding a diploid zygote that develops into a diploid female. Sexual reproduction occurs primarily when living conditions are unfavorable. Most structures in rotifers are syncytial ("a mulitnucleate mass of protoplasm not divided into separate cells," or "a multinucleated cell") and show eutely (here, "constant or near-constant number of nuclei"). Phylum Nematoda. See text pages 894 - 895. Name from the Greek for "thread". This phylum consists of the round worms. There are about 12,000 named species but the true number probably is 10 - 100 times this! These animals are triploblastic, bilaterally symmetrical, unsegmented pseudocoelomates. They are vermiform, or wormlike. In cross-section, they are round, and covered by a layered cuticle (remember this cuticle !!). Probably due to this cuticle, juveniles in this phylum grow by molting. The gut is complete. They have a unique excretory system but they lack special circulatory or gas-exchange structures. The body has only longitudinal muscle fibers. The sexes are separate. Nematodes can be incredibly common, widespread, and of great medical and economic importance. They are parasites of humans and our crops. They can live pretty much anywhere. Nematodes can be free living or important parasites of our crops, or of humans and other animals. They have become very important in development studies, especially the species Caenorhabditis elegans, presumably due to its small size and constancy of cell number (eutely - 959 cells in C. elegans). Phylum Annelida. See text pages 906 - 909. Name means "ringed", from the Greek annulatus. This phylum consists of earthworms, leeches, and various marine worms given many different names (e.g., sand worms, tube worms). There are about 12,000 - 15,000 species. Animals in this phylum are triploblastic, bilaterally symmetrical, segmented coelomates. Development is typically protostomous. They have a complete circulatory system, and a well-developed nervous system. Typically, each segment has paired epidermal "bristles" (setae or chaetae). Most are marine but they are successful occupants of almost anywhere sufficient water is available. They can be free living, parasitic, mutualistic, or commensalistic. Major advances of this phylum include the true coelom, segmentation, both longitundinal and circular muscles, a closed circulatory system and, for most, a more advanced excretory system (metanephridia). There are three main classes of Annelids Phylum Arthropoda. Name means "jointed feet". This phylum consists of spiders, ticks, mites, insects, lobsters, crabs, and shrimp, and is the largest of all the phyla. So far, over 1 million species have been named, and it is likely that the true number out there is 10 - 100 times greater. These animals are triploblastic, bilaterally symmetrical, segmented, protostome coelomates. The coelom is generally reduced to portions of the reproductive and excretory systems. They have an open circulatory system. The most notable advancement of this phylum is a rigid exoskeleton. It has major implications in these organisms' locomotion, flexibility, circulatory systems, gas exchange systems, and growth. It also was partially responsible for the ability of the arthropods to move on to land. There are several major groupings of arthropods: Phylum Mollusca. See text pages 900 - 905. Name means "soft". This phylum consists of snails, slugs, bivalves, chitons, squids, octopus, and many others. About 110,000 species All molluscs have a similar body plan: Molluscs are bilaterally symmetrical, or secondarily asymmetrical. They are coelomates, but the coelom generally has been greatly reduced; the main body cavity is a hemocoel. Development is typically protostomous. The gut is complete with marked regional specialization. Large, complex, metanephridia (excretion). Many molluscan life cycles include a trochophore larva. This stage also is characteristic of annelids. There are several major classes of molluscs: Phylum Echinodermata. Name means "spiny skin" This phylum consists of sea stars, brittle stars, sea urchins, and sea cucumbers. Echinoderms are mostly sessile or very slow moving animals. As adults, they are radially symmetrical, but in the larval stage, they are bilaterally symmetrical. They are considered deuterostomes. Echinoderms are unique in that they have a water vascular system composed of a system of fluid-filled canals. These canals branch into the tube feet, which function in feeding, locomotion, and gas exchange. There are six major classes of echinoderms: Phylum Chordata. Name means "the chordates", i.e., these animals have a notochord at some stage in their lifecycle. This phylum consists of tunicates, lancelets, and the vertebrates. There are four major features that characterize the phylum Chordata. Chordates have a segmented body plan, at least in development. This segmentation evolved independently from the segmentation of annelids. Three subphyla make up the phylum Chordata: Formally, the phyla Urochordata and Cephalochordata are considered invertebrates. Subphylum Vertebrata. Vertebrata refers to the presence of vertebrae and a vertebral column. This subphylum includes most of the animals with which most people are familiar. The notochord generally is replaced by the cranium &amp; vertebral column in adults. Neural Crest Cells. Later in development, these give rise to many cells of the body, including some cartilage cells, pigment cells, neurons &amp; glial cells of the peripheral nervous systems, much of the cranium, and some of the cells of the endocrine system. Some scientists would like to classify the neural crest as the fourth germ layer. Neural crest cells come from the dorsal edge of the neural plate, thus ectoderm. 

Key Terms. synapomorphy Introduction. What makes an animal an animal? If animals are a monophyletic taxon, then animals should be able to be defined by synapomorphies, (shared, derived characteristics). Ideally, we would NOT define this or any taxon using symplesiomorphies (shared ancestral or primitive characteristics) or homoplastic characters (the independent evolution of similarity, or "convergent evolution"). See pages 654 - 656 and Fig. 32.6 in your text to review these concepts. As you consider the characteristics listed below, ask yourself whether or not each is a synapomorphy. Characteristics of an Animal. Animals are multicellular heterotrophic eukaryotes Animals share unique characteristics Animals share certain reproductive characteristics Other commonly used definitions or characterizations What kinds of animals are there? "This text is based on notes very generously donated by Ralph Gibson, Ph.D. of the Cleveland State University." 

The Evolutionary Tree in Animals. There are many competing hypotheses for the form of the evolutionary tree of animals. A traditional hypothesis is that the tree resembles a tuning fork: it has a short base and two main branches. However, there is recent molecular evidence that challenges part of this traditional hypothesis. Under the tuning fork model, the "base" of the tree includes structurally simple animals like sponges, corals, and their relatives. One main branch includes arthropods, molluscs, annelids, and nematodes. This branch, or a large part of it, usually is called the protostomes. The second main branch includes vertebrates (phylum Chordata), and starfish, sea urchins, and their relatives (phylum Echinodermata). This branch usually is called the deuterostomes. Flatworms (phylum Platyhelminthes), which include free living planarians as well as parasitic flukes and tapeworms may be placed very low on the protostome branch, or high on the trunk just below the protostome - deuterostome branching. Support for the "Tuning Fork" Model. The features of animals that have been interpreted as suggesting a tuning fork model are extremely basic characteristics of body organization and early embryonic development. Body Symmetry. Asymmetry. Lack of any symmetry. Many sponges are asymmetric. Bilateral Symmetry. There is only one plane of symmetry, and it is anterior-to-posterior, dorsal-to-ventral, through the midline. Characteristic of most protostomes and the higher deuterostomes. Presence of True Tissues. Tissues are defined as an integrated group of cells that share a common structure and a common function (for example, nervous tissue or muscle tissue). Sponges are described as lacking true tissues. True tissues are present in Cnidaria, flatworms, and all higher animals. Number of Embryonic Germ Layers. Germ layers are defined as the basic tissue layers in the early embryo which give rise developmentally to the organs and tissues of the adult (e.g., ectoderm, mesoderm, endoderm). This is a concept that is applied "only" to organisms considered to have true tissues. Two Germ Layers. Organisms with two germ layers are said to be diploblastic. This is characteristic of Cnidarians. Three Germ Layers. Such organisms are said to be triploblastic. This is characteristic of flatworms and all higher organisms. Four Germ Layers (?). Some developmental biologists consider the neural crest tissue of vertebrate embryos to be a fourth germ layer. Nature of the Main Body Cavity. Most triploblastic animals have a fluid-filled space somewhere between the body wall and the gut. Such a cavity can provide numerous functional advantages. For example, peristalsis of the gut need not affect the body wall, and movements of the body wall during locomotion need not distort the internal organs. We will consider three conditions with respect to the body cavity: Acoelomates. These animals lack an enclosed body cavity; the only "body cavity" is the lumen of the digestive tube. The space between the gut and the body wall is filled with a more or less solid mass of mesodermal tissue. The major example of this is the phylum Platyhelminthes, the flatworms. Minor examples: Phylum Nemertea (Rhynchocoela) and Phylum Gnathostomulida (not responsible for these "minor" examples). Pseudocoelomates. Pseudocoelomates have a fluid-filled cavity between the body wall and the gut, but it does not form within mesoderm, nor does it end up fully enclosed by mesoderm. This cavity often is interpreted as being a developmental remnant of the blastocoel, the fluid-filled cavity of the blastua stage of the embryo. To distinguish it from the next grade, this type of cavity is called a pseudocoelom. Major examples of pseudocoelomates include the phyla Nematoda (round worms) and Rotifera (rotifers). Other phyla listed in the table that are considered to be pseudocoelomates are flagged with an asterisk. Note that recent molecular data in particular have challenged the "naturalness" of the pseudocoelomates as a possible taxon. Coelomates. Coelomate animals also have a fluid-filled cavity between the body wall and the gut. In this grade, however, the cavity is completely enclosed by mesoderm. Major examples of coelomates include molluscs, arthropods, echinoderms, and chordates. The protostome-deuterostome distinction. The distinction is based on several fundamental characteristics of early development. Characteristic: Determinate vs. indeterminate cleavage Spiral vs. radial cleavage Fate of the blastopore Source of mesoderm Formation of coelom Phyla Echinodermata, Hemichordata, Chordata "This text is based on notes very generously donated by Ralph Gibson, Ph.D. of the Cleveland State University." 

__NOEDITSECTION__ Structured Query Language (SQL) is a widely-used programming language for working with relational databases. The name of the language is generally pronounced as the three letters of its abbreviation or, in some people's usage, as . This Wikibook provides a short description of SQL, its origins, basic concepts and components, and many examples. The book follows the specifications of the SQL:2011 standard developed by a common committee of ISO and IEC. Their publications are not freely available but can be ordered online. Or you may want to refer to a working draft that you can download from Whitemarsh Information Systems Corporation. 

Scope. This work is aimed at people already familiar with using the Internet, who want to know how and why it works. When we say technology we don't just mean the software and hardware, but also the human components which are an integral part of the overall system of the Internet. 

Spanish slang is more localized than English slang and sometimes people from one Spanish-speaking country get confused talking to people from other Spanish-speaking countries. 

" The World History Project" Welcome to the World History Project. This organization is dedicated to making a free, open-content, standardized textbook on World History based on the AP World History Standard. The goal is to create a standard of quality which will suffice for a secondary and post-secondary environment. The World History Project is the "brains" behind the organization. We are a set of regular contributors who organize and give the major guidance to the World History page. We welcome contribution of any who wish to help (whether as part of the World History Project or no), as well as collaboration with other projects - contact us at our main discussion page or here at our Authors page. Standards | Our Golden Rule | The Authors | Prologue - Other.  __NOEDITSECTION__ 

This page is for an elementary explanation of the origins and nature of the religions. It is not meant to be an in depth study by any means. Animism. Animism (from Latin animus, -i "soul, life") is the worldview that non-human entities (animals, plants, and inanimate objects or phenomena), possess a spiritual essence. Baha'i Faith. The "Bahá'í Faith" is the youngest of the world's independent religions. Its founder, (1817–1892), is regarded by Bahá'ís as the most recent in the line of that stretches back beyond recorded time and that includes Abraham, Moses, Buddha, Zoroaster, Christ and Muhammad. The central theme of Bahá'u'lláh's message is that humanity is one single race and that the day has come for its unification in one global society. God, Bahá'u'lláh said, has set in motion historical forces that are breaking down traditional barriers of race, class, creed, and nation and that will, in time, give birth to a universal civilization. The principal challenge facing the peoples of the earth is to accept the fact of their oneness and to assist the processes of unification. Wikipedia article: Buddhism. Buddhism is a religion based largely around the teachings of the Siddhārtha Gautama (although his exact status is still controversial and changes by sect) and is now the central religion of most of Southeast Asia, Mongolia, Sri Lanka, and parts of China, Nepal, India, and a small part of Russia. It is also a major religion in Korea and Japan and has a growing influence in the west. Gautama was born in ancient Nepal around the 6th century BC, son of a King and relieved of all tasks. One day, while being brought around by his charioteer, he saw the four passing sights (an old man, a sickly man, a decaying corpse, and a holy man), which brought him to the realization that birth, old age, sickness, and death happen to all people over countless lives. He left his wife, children, rank, and his entire life to solve that problem. Gautama tried everything to achieve inner peace (he nearly killed himself numerous times) but found nothing that worked. He then tried sitting peacefully under a Bohdi tree and meditating. This proved very successful, and he soon achieved the inner peace he wanted. He then traveled the lands, preaching his new faith. (Postscript-Buddhism is largely based on Jainism, and shares many beliefs with Jainism.) Within the context of postclassical China, dominations such as pure land and Zen Buddhism appealed to both aristocratic elites and the mass peasantry. A commonality of religion in the global context: the induction of fervent belief system in the presence of societal corruption and lack of intellectual synthesis. Modern Buddhism still follows the ideals of Gautama-peace, kindness to man, and love of nature (including vegetarianism). There are three modern sects of Buddhism - Theravada, Mahāyāna, and Vajrayāna (practiced in southeast Asia, East Asia, and scattered parts of Asia, respectively). For information on the sects of Buddhism, or some tenents of Buddhism, go to: Buddhism Theravada Mahayana Vajrayana Christianity. Christianity Confucianism. Confucianism Hinduism. "Hinduism" is a term coined to designate the traditional socio-religious systems of the people of India. This term does not appear in any of the sacred literature of India. Hindus refer to their religion as "Sanatana Dharma" which loosely translated means “The Eternal Path”. "Sanatana" means "eternal", "perpetual" or "sustained". "Dharma" means any method by which one sees reality for what it is, and that by which one is drawn closer to the Absolute Truth and Ultimate Reality — it is the "Philosophia Perenis". In a context of world history, the Hindu emphasis placed upon social divisions as ample means for a productive society led to the highly stratified caste system in which birth and socio-economic position determined semi-permanent placement. There are two world religions which have formed the cultural and ethical basis of the world as we know it. Both have an unbroken history going back thousands of years. Judaism with a 5000 year old tradition is the mother of the western civilisation through its offshoot Christianity. Hinduism is the older of the two with a literature going back to the beginning of recorded history. Hindu civilisation originated in the Gangetic and Indus valleys and from there spread out over the entire region of southeast Asia. Its offshoot — Buddhism, shaped and molded the civilizations of Japan, China, Tibet and the rest of Asia. There is evidence to suggest that the Ancient pre-Biblical kingdom of the Mittani in Asia minor was ruled by Kings with Hindu/Sanskrit names. The Hittites were an Indo-European people and according to some sources are said to have originated in the Gangetic Basin of India. Hindu philosophy/theology influenced the ancient Greeks since the time Alexander the Great conquered parts of north India. A remarkable similarity has also been demonstrated between the religion and mythology of the ancient Scandinavian people and that of the people of India. The ancient civilizations such as the Roman, Greek, Egyptian, Sumerian, Babylonian, Mayan, Aztec, and Inca have all passed away. Even the Jewish culture has undergone many radical changes since its inception 5000 years ago – yet the Hindu civilisation continues as a vibrant and living vector, and has remained virtually unchanged for over 6000 years. Today, Hindu communities are to be found in almost every country on earth. Hinduism Islam. Islam Jainism. Jainism Judaism. Judaism Sikhism. Sikhism Shinto. Shinto Taoism. Taoism/Daoism 

10.2 Students compare and contrast the Glorious Revolution of England, the American Revolution, and the French Revolution and their enduring effects worldwide on the political expectations for self-government and individual liberty.. 1. Compare the major ideas of philosophers and their effects on the democratic revolutions in England, the United States, France, and Latin America (e.g., John Locke, Charles-Louis Montesquieu, Jean-Jacques Rousseau, Simón Bolívar, Thomas Jefferson, James Madison). 2. List the principles of the Magna Carta, the English Bill of Rights (1689), the American Declaration of Independence (1776), the French Declaration of the Rights of Man and the Citizen (1789), and the U.S. Bill of Rights (1791). 3. Understand the unique character of the American Revolution, its spread to other parts of the world, and its continuing significance to other nations. 4. Explain how the ideology of the French Revolution led France to develop from constitutional monarchy to democratic despotism to the Napoleonic empire. 5. Discuss how nationalism spread across Europe with Napoleon but was repressed for a generation under the Congress of Vienna and Concert of Europe until the Revolutions of 1848. 

Objectives. 10.3 Students analyze the effects of the Industrial Revolution in England, France, Germany, Japan, and the United States. 1. Analyze why England was the first country to industrialize. 2. Examine how scientific and technological changes and new forms of energy brought about massive social, economic, and cultural change (e.g., the inventions and discoveries of James Watt, Eli Whitney, Henry Bessemer, Louis Pasteur, Thomas Edison). 3. Describe the growth of population, rural to urban migration, and growth of cities associated with the Industrial Revolution. 4. Trace the evolution of work and labor, including the demise of the slave trade and the effects of immigration, mining and manufacturing, division of labor, and the union movement. 5. Understand the connections among natural resources, entrepreneurship, labor, and capital in an industrial economy. 6. Analyze the emergence of capitalism as a dominant economic pattern and the responses to it, including Utopianism, Social Democracy, Socialism, and Communism. 7. Describe the emergence of Romanticism in art and literature (e.g., the poetry of William Blake and William Wordsworth), social criticism (e.g., the novels of Charles Dickens), and the move away from Classicism in Europe. The Industrial Revolution. In 1750, most people in Europe lived on small farms and produced most of their needs by hand. By the middle of the 19th century, many people lived in cities and most of their needs were produced by complex machines using steam power. The Industrial Revolution began in Great Britain and spread to Belgium, France, Germany, the United States and Japan. It was a fundamental change in the way goods were produced, and altered the way people lived. The Industrial Revolution is a major turning point in world history. Why Great Britain? Great Britain became the focus of the Industrial Revolution for a variety of reasons: the start of the Agrarian Revolution, an abundance of natural resources, available capital, and the political will to support innovation. The Agrarian Revolution was a change in farming methods that allowed for a greater production of food. This revolution was fueled by the use of new farming technology such as the seed drill and improved fertilizers. The results of this revolution in farming was a population explosion due to the higher availability of food. Also, the Enclosure Movement, which was the consolidation of many small farms into one large farm, left many people jobless and homeless. These people provided the workforce of the Industrial Revolution. Great Britain's geography provides them with an abundance of the natural resources needed for industrialization, such as iron ore and coal. Britain also had access to many navigable rivers and natural harbors which provided for the easy movement of goods both within the country, and overseas. The British overseas empire provided them with a strong economy, this produced the capital (money) needed to build railroads, factories, and mines. Politically, British entrepreneurs enjoyed a high degree of freedom from state control, compared to their counterparts in France, Russia and other parts of Europe. A relatively fair court system existed to enforce contracts and settle disputes among capital owners. These factors may have allowed new technologies and energy resources to take root and flourish. Britain experienced a revolution in energy use as they switched from animal power, to water power, to steam power in a few short years. The steam engine was the power source of the Industrial Revolution. Effects. Philosophy. The philosophy of Communism appeared as a reaction to the condition of the Working Class in industrial society. Karl Marx wrote in The Communist Manifesto (1848) that all of human history is based on the conflict between the bourgeoisie (those who control the means of production) and the proletariat (working class). He predicted that the proletariat would rise up in a violent revolution to overthrow the bourgeoisie and create a society with an equal distribution of goods and services. This socialist theory would form the basis for the Bolshevik, Chinese, and Cuban Revolutions in the 20th Century. The United States had a very strong reaction to him. Imperialism. Due to the need for raw materials and new markets, the industrialized nations took control of Africa, India, South East Asia, and others. Imperialism had a negative effect on most of these cultures, and did not completely end until after World War II. Most of the benefits of imperialism accrued to the European nations. The Industrial Revolution was a major turning point in world history as it resulted in a complete change in society on all levels. Effects of the Industrial Revolutions were long reaching, and influenced many other cultures both positively and negatively. 



Before the Great War. Factors Leading to War. Constant colonial tensions among the great powers had given rise to the possibility of a great war between the major European powers. For almost a hundred years, since the fall of Napoleon, a remarkable series of events kept the relative peace. But as 1914 approached, tensions began to rise in a number of countries, and key issues began to take their toll. The Changing of Imperialism. By 1914, Imperialism had begun to come at last to a turning point. The massive land grabs that had divided up Africa, given Britain such a huge empire, and led to the collapse of the Chinese Empire had finally run their course. There simply were very few places left on Earth containing massive amounts of land not already claimed by a western power. This was especially bad news to nations that were latecomers to the game of Imperialism, Germany and, to a lesser extent, Russia. Wilhelm II did not wish for Germany to miss out on the benefits of a colonial empire. Germany already had some territories in Africa, and after the Berlin Conference of 1885, they were major policy makers on the continent. However, Wilhelm II continually pushed for a larger role for Germany in Africa, leading ultimately to the Morocco Affair, and a heightening of tension. Turmoil in Russia. Russia was also undergoing immense changes. After three hundred years of Romanov rule, Russia was finally making a transition to a modern nation. However, the Romanov tsar, Nicholas II of Russia , was still clinging onto a good measure of authoritarian power. The struggle between Nicholas II and reform-minded progressives would eventually lead to the drama of the Russian Revolution. To make a bad situation worse, Russia was also in a bad geographic situation. Despite being the largest European power in terms of land area of their home nation, Russia was mostly poor and un-industrialized. Much of the nation was underdeveloped, and only slowly inching its way toward a more profitable existence. At the same time, Russia was dealing with an unfavorable political situation, already reeling from the loss of the Russo-Japanese War, as well as attempting to stay afoot of the fast-paced diplomacy of the west. The Sick Man of Europe. For several hundred years, the Ottoman Empire had been slowly collapsing under its own weight, watching helplessly as first one province and then another had broken away, or been stolen. Recently, in the Crimean War, the Ottomans had been forced to rely upon the aid of Britain and France to sustain itself in conflict against the Russians. The Empire had previously been called the "Sick Man of Europe", but prior to the war it was called the "Dying Man of Europe". Only sparse pieces of the Balkans remained in Ottoman control by 1914, and they were being encroached on from many fronts. The Austrians in the north clearly wanted to pacify provinces in the region, the Russians in the east wanted Istanbul itself to guarantee themselves a safe passage through to the Mediterranean from their Black Sea ports, and the people who lived in the Balkans were beginning to experience their own internal unrest. By 1914 it looked like conflict could erupt at any moment in the Balkans. Naval build up and the end of Splendid Isolation. In order to demonstrate their military capability, Germany embarked on a plan to build up her navy, constructing the High Seas Fleet, equipping it with the latest in military technology. This not only contributed to a rise in Germany's military power, but it also seriously alarmed Great Britain. Britain had long subscribed to the theory of "Splendid Isolation", under which Britain attempted to hold itself aloft from affair on the continent. Protected by the Royal Navy, the most powerful navy in Europe, Britain had remained inviolate to foreign powers for centuries; even Napoleon was unable to land an army on her shores. Under Splendid Isolation, Britain chose to abstain from permanent relationships with European powers and to continually work to maintain a balance between them, all the while protected from continental strife by her powerful navy. The construction of the German High Seas Fleet pushed Britain into a fury of shipbuilding, in an attempt to stay ahead of German production. The knowledge that they could no longer stand by in the distance and remain safe from European struggle, and the continuing threat created by German armament programs, propelled the British out of their self-imposed isolation. It also led to a souring of relations between both Britain and Germany. Assassination of Archduke Franz Ferdinand. On June 28, 1914, at approximately 11:00 am, Franz Ferdinand and his wife were killed in Sarajevo, the capital of the Austro-Hungarian province of Bosnia and Herzegovina, by Gavrilo Princip, a member of Young Bosnia and one of several (seven) assassins organized by The Black Hand (Crna Ruka). The event, known as the Assassination in Sarajevo, was one of the main triggers of World War I. War Breaks Out. Officially, the First World War began on June 28, 1914 in the Bosnian city of Sarajevo, with the assassination of Archduke Franz Ferdinand of the Austro-Hungarian Empire by Gavrilo Principe, a member of a secret Serbian society, The Black Hand, hoping to unite the Slavic speaking lands of the Austro-Hungarian Empire with the nearby Kingdom of Serbia. In retaliation, Austria, with the support of Germany, presented Serbia with a list of demands that would severely threaten the Kingdom's autonomy. The Austrian leaders had been waiting for this moment for several years, because they feared just the sort of Serbian groups that were now operating in Bosnia, and had been waiting for this chance to weaken Serbia for several years. After consultations with the Russians, Serbia rejected the demands, leading to war with Austria. Russia responded with a general mobilization, followed by Austria's ally Germany and Russia's ally France. The first stage of the war quickly escalated, with a general war breaking out between an alliance of Germany and Austria against Russia, France, and Serbia. The war presented the Germans with a special problem. Their own armies were scattered on the country's eastern and western borders against France to the west and Russia to the east, but the armies of France and Russia were concentrated against the Germans on their respective borders. In order to make the best out of this two front situation, the German generals began to carry out the Schlieffen Plan, developed by a German general in case of war with both France and Russia. The plan called for a quick and massive strike against the French while the massive Russian army was still mobilizing, hoping that their Austrian allies could stop the Russian advance long enough for them to finish off the French and move their troops by rail to defend the east. Unfortunately, the quickest way to strike the French capital, Paris, and knock the French out of the war involved a broad offensive through neutral Belgium. Great Britain, a nominal ally of France, promised to intervene in case Germany invaded Belgium, and declared war shortly after German forces crossed the border. The Germans initially expected the Belgians to put up very little resistance, but this proved to be incorrect and the longer than expected time required to defeat the Belgians allowed the French to regroup. While the Germans were pushing deep into Belgium, Serbian forces crossed the Danube in a daring attack against the Austrians. This offensive, like many others after it, failed, and the war zone along the Balkans would change very little during the first year. Meanwhile, the massive German offensive bearing down on France was beginning to run out of steam. First, the Russians scored several early victories against the Austrians and moved through Poland (then a Russian territory) on their way to attack Prussia and the vulnerable German capital, Berlin. The worried German generals withdrew several divisions to the eastern front, allowing the French to regroup and defeat the German lead elements at the Battle of the Marne, slightly more than 15 miles outside Paris. Though the German armies were stopped by the French, the Russian armies bearing down on Prussia were similarly defeated by German forces under Field Marshall Hindenburg at the Masurian Lakes, and then at Tannenberg. Though the Russian army was much larger than the German army, the Russian commander's decision to split his forces in a two-prong attack contributed to his defeat, and allowed the Germans to break apart the offensive piecemeal. Stalemate and Trench Warfare. Following these defeats, the war on both fronts bogged down. In the west, both French and British forces and the Germans built a massive series of trenches from the Swiss border to the North Sea to defend their positions. This period of "Trench Warfare" saw very little change in the battle lines, and several uses of poison gas by the Germans. With no major breakthroughs on either fronts, and furthermore the people in Russia were killed. German and British leaders looked to the Mediterranean to break the stalemate. This put the two remaining neutral powers, Italy and the Ottoman Empire (later Turkey) in a tough position. Before the war, the Italians had agreed to support Germany and Austria in case of war. However, British and French diplomats promised the Italians the Austrian territories adjacent to Venice if Italy declared war, which they did after the prodding of many Italian radicals, among them Mussolini. The Ottoman Empire also had to make a tough choice, because Britain was the Empire's traditional protector and Russia was the Empire's traditional enemy. Germany's offer of two battleships and Middle Eastern territories allowed the Turks to strike out at both simultaneously, beginning with a naval strike against the Russian naval base at Odessa. The Ottoman entry into the war encouraged neighboring Bulgaria to also enter the war in an alliance with Germany, leading to the rapid defeat and occupation of Serbia by the Austrians. Austria played an important part in the war and still is today a leading military nation. In 1915 the Ottomans launched an attack into Egypt in an attempt to capture the Suez Canal. The assault failed and the Ottomans retreated back into their empire. Once again in 1916 the Ottomans attempted to capture the canal, but this once again failed. This attack convinced the British to push their defence of the Canal further out, into the Sinai, and so starting in October, the British under Lieutenant General Sir Charles Dobell began operations into the Sinai desert and on to the border of Palestine. Initial efforts were limited to building a railway and a waterline across the Sinai. After several months building up supplies and troops, the British were ready for an attack. The first battle was the capture of Magdhaba on December 23 1916. This was a success, the fort was captured. In 1916, the Central Powers invaded Romania, and forced them to sign an unwilling peace. America Enters the War. By 1916, the two groups had been fighting for two years with neither side making significant gains. Food shortages had appeared in Germany (under blockade by Britain) and Russia, leading to much discontent. Responding to the British blockade, the Germans launched a naval offensive of their own, using submarines known as U-Boats to prey on allied shipping. The Lusitania. A British-American tourist cruise liner that was intentionally sunk by a German U-Boat. However, there was some ammunition onboard in the cargo area, and this was the German reason for sinking. The Zimmerman Telegram. A telegram from Germany to Mexico that encouraged Mexico to attack the United States. This telegram was provided to the United States by Great Britain and there has been some speculation on whether it was forged by Great Britain in order to get the previously neutral United States to join World War I Defeat of the Central Powers. Germany miscalculated by launching the u-boat campaign, which intended to win the war before the entry of USA. This miscalculation was a military suicide on the part of Germany as it led to the entry of USA. America with her vast resources entered the war on the side of allied forces leading them to a sweet victory against Germany. 

Quick Quiz Benito Mussolini and Fascism in Italy (1922-1939). Benito Mussolini, born into a poor blacksmith's family, was so named by his radically socialist father (his mother was a devout Catholic schoolteacher) after the executioner of a Mexican emperor. Shortly after becoming qualified as a teacher, Mussolini taught in a small school. Mussolini was a far-left socialist and advocated a violent revolution to overthrow the parliamentary monarchy within Italy and denounced nationalism. When World War I broke out in 1914, however, he broke with his party comrades when he celebrated the entry of his nation into the war – even though he had dodged the draft. Throughout the Great War, he fought earnestly to keep Italy involved, and, financed by large arms manufacturers and the British and French governments, operated a small, pro-war newspaper. When the war came to an end in 1919, Mussolini was quick to recognize the dissatisfaction of many of the homebound soldiers and countrymen concerning the Treaty of Versailles. In an effort to persuade Italy to enter the war on their side, the Allied Powers promised Italy significant territorial gains at the expense of the Austro-Hungarian Empire. The final settlement, however, was less favorable to Italian interests than that originally promised, and resulted in widespread malcontent regarding the post-war government. In March 1919, Mussolini created a radically nationalist and anti-communist party – Fasci Italiani di Combattimento. Mussolini, who loved the splendor and extravagance of Ancient Rome, adopted a Roman symbol of authority, the "fascio" (an axe wrapped in whipping rods) for his group of devotees. As inflation and economic decline spread throughout Europe and Italy following the war, factory workers began to go on strikes in northern Italy. In 1920, Mussolini’s group’s numbers were bolstered by ex-soldiers willing to break up these strikes. Mussolini marched 50,000 Fascist supporters (known as Blackshirts for their attire) in squads against the strikers and left-wing newspapers. The Blackshirts garnered their support from the financial contributions of industrialists and large landowners, who shared their anti-communist sentiments, but also believed that they could control the excesses of the Fascist party. The police often refused to stop the squads, allowing the Blackshirts freedom to inflict whatever damage they wished. The widespread destabilization of the previous social orders throughout Europe due to economic uncertainty in the aftermath of the war and the successful establishment of the Soviet Union as a socialist state led many to believe that democracy was weak and ineffectual, while monarchy was discredited as an oppressive and unresponsive system. A command economy was thought to be a progressive and scientific method of social organization. Fascism incorporated the futuristic and populist elements of Communist ideology, but also identified itself strongly with the nationalism that had created the modern European nation-states in the late 19th Century. Despite growing popularity and the introduction of proportional representation in the Parliament, Mussolini's party fared poorly at the polls, winning no seats in 1920 and only 35 in May 1921 (7% of the vote). The internal political situation however, swung in Mussolini's favor. The birth of the Communist Party of Italy, openly allied with Lenin's Soviet regime in Moscow, polarized Italian politics. Proportional representation caused stagnation in government, until a weak coalition finally came into place in February 1922. In October of that year, as Mussolini was giving one of his soon-to-be characteristic speeches from atop a balcony, he suddenly cried, “To Rome! To Rome!” The crowd of supporters, much to Mussolini’s surprise, echoed his cry. Blackshirts to the number of 40,000 organized to march on the capital. Mussolini, however, went into hiding, afraid of the impending collapse of his movement. When it became clear that the army would not oppose his action, however, Mussolini moved decisively. As Blackshirts began to occupy key posts in Rome on 27-28 October, the king, Victor Emmanuel III, to the chagrin of his elected cabinet, appointed Mussolini as Prime Minister so long as Mussolini halted the advance to Rome. Mussolini agreed. His word was not to be trusted, however, as he soon after marched on the city anyway, creating an incredible propaganda success for the Fascists. Over the next few years, he led a slow-motion coup d'état. By 1926, he had become the undisputed totalitarian dictator (“Il Duce,” or “the leader”) of Italy. Mussolini's regime embarked on a campaign of militarization and political maneuvering. First, he began to ready Italy for war. One of Mussolini’s driving ambitions was to restore the hegemony of the Roman Empire in a modern Italy. To that end, he encouraged couples to have as many children as possible as he organized large-scale expansions of the agricultural sector to feed them. He extended an olive branch to the Catholic Church by way of the Lateran Accords, which recognized papal authority over the Vatican and declared Catholicism the official religion of all of Italy, ending hundreds of years of estrangement. Mussolini seemed to have been a victim of his own propaganda as, in 1935, he deemed his newly-formed army strong enough to invade Ethiopia. The Italian army invaded from Italian-held Eritrea. The underestimation of the enemy proved fatal for thousands of ill-prepared Italians as the army met face-to-face with the “Lion of the Desert,” Omar Mukhtar (whose death by hanging at the hands of the Italians ended his twenty-year resistance). The poorly-armed Ethiopians were eventually defeated primarily due to terror tactics, such as poison gas and terror bombings. Mussolini’s sense of superiority did not seem hurt by his army's poor preparation, as, over the next few years, he became a close ally of Hitler's. In 1939, Mussolini signed the “Pact of Steel,” creating a formal alliance between Italy and Nazi Germany. Mussolini, while publicly effervescent, did not have the universal power and control enjoyed by Adolf Hitler and Joseph Stalin, a lack of gravitas which would later cost him his life. Nevertheless, Mussolini rode at the helm of the 20th century dictatorship, invented the terms Fascism and Totalitarianism, and pioneered the use of propaganda to control the masses in newspapers, posters, radio, and in movie cinemas. The Weimar Republic (1918-1933). Following the complete collapse of Germany's armed forces throughout the waning months of 1918, German generals and politicians desperately sought to surrender. The Allied Powers, however, would not negotiate with the autocratic Kaiser Wilhelm II, and insisted upon Germany to adopt a democratic government. In this disarray, Germany quickly fell along the slippery slope to revolution, appearing as though it might go the same direction as Russia – Marxist. After Kaiser Wilhelm II abdicated the throne on November 7, a new republic was declared, and a National Assembly convened at Weimar to circumvent the unrest in Berlin. The hastily-formed republican government took its name from the host city and surrendered. The Allied Powers, in turn, forced upon the defeated nation extremely harsh and punitive terms through the Treaty of Versailles. Clause 231 of the treaty, the so-called "war-guilt" clause, called for Germany to accept total and sole blame for the Great War (while, arguably, they merely joined) and to pay reparations to the “victimized” Allied Powers. Finally, it was established that Germany was to disband its air force permanently, and to have no more than 100,000 men in its armed forces. The Rhineland, along the Franco-German border, was demilitarized and put under French jurisdiction. The extremely valuable Saar region, home to most of Germany’s factories, was made autonomous. These harsh restrictions gave the fledgling Weimar Republic unwarranted disrespect. Many Germans had opposed the treaty, and it created large amounts of resentment within Germany. (Retrospectively, the Treaty of Versailles is seen as one of the most fundamental causes of the rise of Adolf Hitler and World War II.) The newfound lack of industry compounded the reparations Germany was forced to pay. To assuage its monetary woes, The Weimar Republic began to print paper money at exorbitant rates – rates so high, that, by 1923, the American Dollar was worth 4.2 trillion German Marks. Amidst the chaos of disappearing life savings and a tumultuous economy, Gustav Stresemann came to the forefront of Weimar Politik. Under his leadership, the Weimar Republic managed to regain marked stability in the period of 1923-1929. Peoples discontent about the Weimar government increased day by day. Hyperinflation was corrected, but Stresemann's death in 1929 and the catastrophic worldwide Great Depression the same year brought about the death of the Weimar Republic. This untimely fall led to the empowerment of a man who would vault Germany to an unprecedented world power, who would pursue the elimination of “undesirables” such as Jews and homosexuals, who would begin the second world war of the 20th Century. Hitler, who had been subverting many of his countrymen during the economically tumultuous 1920s, took advantage of the Weimar Republic’s fall. The Rise of Adolf Hitler in Germany (1914-1939). In August 1914, as the world took the fatal plunge into World War I, an unknown and unimportant young Austrian national named Adolf Hitler enlisted in the German Army. Born on April 20, 1889 into a troubled and strict Austrian family, Hitler was a failed artist and an ardent German nationalist (Austrians are ethnically German and indistinguishable from their cousins). His anti-semitic views already in place from his early life as a vagrant (he dropped out of high school and was refused admission to a Vienna art school), Hitler was eager to serve his adopted homeland. He had an exemplary record of service and received the prestigious Iron Cross, both First and Second Class, and also achieved the undistinguished rank of Corporal. Shocked and deeply angered by the German defeat in 1918, he personally put the sole blame on the so-called "November politicians" (referring to those who formed the Weimar Republic). He also put blame on the Jews for the downfall of Germany. After the war, Hitler remained in the army and after receiving intelligence and oratory training, became an intelligence official tasked with infiltrating political parties and reporting to his superiors on their activities. In March 1919, he was instructed to sit in on a meeting of the small nationalist German Worker's Party. He joined the party in September, and upon his discharge from the army in 1920, soon became the leader of the party which changed its name to the German National Socialist Worker's Party (NSDAP or Nazi for short, from its German name "Nationalsozialistische Deutsche Arbeiterpartei"). Over the next few years, Hitler's oratorical skills allowed the party to expand. It soon had its own private armed forces, known as the SA led by Ernst Rohm. Another important admirer was Erich Lundendorff, a Field Marshall from the First World War, whose help proved invaluable in setting up the Beer Hall Putsch. The Beer Hall Putsch and Mein Kampf (1923-1925). On November 8, 1923, Adolf Hitler and a group of SA raided a beer hall in Munich where the three most powerful politicians in Bavaria were giving speeches. Taking the men hostage, Hitler threatened them with death (and his own suicide) if they did not side with his intention to overturn Bavaria's government and then to march on Berlin. The men agreed (with little other choice). Hitler then made the colossal error of leaving the hall. He left Marshall Lundendorff in command, who upon the assurances of the three politicians that they only wished to return home to their families and would continue to support Hitler, allowed them to leave the hall. The men quickly denounced Hitler and mobilized the government's resistance to his "revolution". Adolf Hitler was enraged. He decided to march his SA the next morning against the Bavarian government. However, army regulars were already at the War Ministry when Hitler arrived and the rebellion was quickly scattered. Hitler was arrested and tried. He spoke so forcefully at his trial however, that the head judge had to harass the other two judges into even convicting him at all. He received a five year sentence. The abortive coup Hitler tried to carry out is referred to as the Beer Hall Putsch. In prison, Hitler dictated the book "Mein Kampf" (My Struggle) to his close friend and confidant, Rudolf Hess. The book was a savage "hymn of hate" denouncing Jews as "parasites" and laying down the foundation for the plan of military conquest Hitler would later attempt. It was all painfully clear: the rearmament of Germany, the invasion of Poland, the invasion of the Soviet Union; Hitler had written down for anyone who wished to read it his plan of action. Unfortunately, few non-Germans read the book, but all too many Germans did. Hitler was released after spending only eight months of his sentence, mostly because the authorities thought he was harmless. He found the Nazi party virtually moribund. In 1925, he formed the "Schutzstaffel" (SS) to be his personal body guard under the leadership of Heinrich Himmler. The Nazi Regime (1933-1939). Through the use of propaganda, Hitler became immensely popular among the German people. To end the depressions, Hitler followed a program of massive public works, including the infamous Autobahn, dams, roads, railroads, and civil improvements. His official announcement of rearmament in 1936 (although it had actually begun much earlier) stimulated the economy further, as it would in the United States during the Second World War. Culture evolved along a strict set of party rules. Men were the heads of work and home; a woman's place was as a cleaner and a mother. The Nazis encouraged large families to literally create men to serve in the army. The Nazis, through their policy of racism, wanted superiority in every sphere of life. When the Olympic Games came to Berlin in 1936, the Germans showed off their athletes in huge stadiums built for the purpose. Adolf Hitler practiced a policy of racial superiority of the Germans, whom he called Aryans, and people were sorted by the correct ethnic "purity". The ideal was the tall, blond, blue-eyed, muscular, and handsome Nordic youth (ironically, Hitler had brown hair). Hitler's regime followed a totalitarian policy; the SS and the secret police, the Gestapo, ruthlessly enforced loyalty to Hitler and rounded up the Nazi's enemies. In 1934, when the army had demanded as the price of its support the dissolution of the SA, Hitler had Ernst Rohm assassinated. Heinrich Himmler became the chief of secret police activities and the mastermind behind the terror. In 1935, the Nazis enacted the Nuremberg Laws, which placed extreme restrictions of Jews and their freedoms as human beings. The economic lives of the Jews were smashed. At this stage however, Hitler was not actively killing Jews but deporting them. The Nazis operated concentration camps at this time to deal primarily with political prisoners. The propaganda machine of the Nazis was similar to that of Stalin in the USSR. However, the Nazis, under propaganda minister Joseph Goebbels, used their propaganda to acquire not only acquiescence to Hitler's schemes, but also to convince the Germans of their policy of racial purity and antisemitism. Goebbels saw to it that, like in the Soviet Union, a picture of the Führer appeared in every building and home, and in many public places. Posters were one of the favorites of the Nazis. They also used the theater extensively to bring in support for the party's goals. Joseph Stalin takes power in the Soviet Union (1924-1934). When Vladimir Lenin died in 1924, he left a power vacuum behind in the wake of his death, centering on the continued use of the New Economic Policy (NEP). The primary contenders for political power were Joseph Stalin and Leon Trotsky. Leon Trotsky was a brilliant politician, and had been Commissar of War during the Civil War. He was a gifted orator and a dedicated Communist, especially to the cause of causing Marxist revolutions internationally, through the use of arms if need be. Ironically, Trotsky had originally been a member of the Menshevik faction of the Russian Social Workers Party, until Lenin, recognizing his genius, had won him over to the Bolshevik camp. Stalin, on the other hand, was a gifted organizer. He was referred to by many of his contemporaries in the party as "Comrade Index-card". However, Trotsky was obviously the more popular choice for the job as the head of the new communist state. Unfortunately for Trotsky, Stalin was also the party's General Secretary. Although primarily a bureaucratic job, the General Secretary actually held the most power in the party because he appointed regional and local party posts in government. Stalin was therefore in a position to appoint those who would support his bid for power. Stalin initially allied himself with the right and center factions of the Communist Party (if any part of a far-left party may be called "right") which supported the continued existence of the NEP. Allying himself with Lev Kamenev and Grigori Zinoviev, he threw his might against Trotsky, who was removed from his post as People's Commissar of War. Stalin now turned against Kamenev and Zinoviev, allying himself with Nicolai Bukharin. Trotsky was expelled from the Communist Party on November 12, 1927, and expelled from the Soviet Union in 1928. He eventually found his way to Mexico, where he was murdered in 1940, probably on Stalin's orders. Now Stalin turned on his allies again, abandoning Bukharin and calling for the abandonment of the NEP. By now, Stalin was the undisputed leading figure of the Communist Party. By the early 1930s, Stalin would truly become the dictator of the Soviet Union. The First Five-Year Plan and Collectivization (1927-1939). At the Fifteenth Congress of the Communist Party, Stalin openly advocated the end of the NEP and introduced a plan for rapidly industrializing the largely rural Soviet Union, remarking that the country was "fifty to one hundred years behind the advanced countries". The government then introduced Gosplan (The State General Planning Commission) which came up with basis for the Five-Year Plan, aimed to turn the country into a major industrial power within five years. The plan set ridiculously high quotas for development. Nonetheless, terrific economic growth was achieved, especially in the areas of coal and iron output. As a result, steel production grew exponentially. However, harsh penalties for not making quotas caused large-scale misrepresentation of growth to occur. Harsh totalitarian measures were introduced. Miners were expected to put in 16 and 18 hour work days, unheard of even the strictest parts of the major capitalist countries. Poor and hazardous working conditions caused countless deaths. Most of the massive industrial complexes constructed for the Five-Year Plan were built by slave-laborers, sentenced for trivial and often completely false crimes. Approximately 3.7 million people were sentenced for counter-revolutionary crimes, approximately 0.6 million were put to death, 0.7 million were expatriated, and 2.7 million were sent to forced labor camps (called the Gulag), often itself a death sentence. As another part of the Five-Year Plan, the government began to forcibly collectivize agriculture (that is, to create large-scale farms where peasants worked the land collectively). The state sought not only an increase in agricultural output, but also to export grain abroad, in order to gain financial capital to buy important technologies for the industrial parts of the Five-Year Plan. By 1936, 90% of the nation's farms had been collectivized. However, this was not done without cost. Peasants almost universally actively opposed collectivization. In the Ukraine, the peasants killed off livestock rather than give it to the authorities. Stalin was so incensed that he allowed a famine to occur which led to the deaths of millions of innocent Ukrainians. As a result, throughout the period of 1924-1953, agricultural output was generally low, not regaining output levels of the period of the NEP until 1940, and rising only marginally in the following years. In addition, Stalin saw fit to deal with richer peasant farmers (known as Kulaks) by deporting them to forced labor in Siberia. In practice however, any person critical of collectivization was deemed a Kulak and summarily deported. It is estimated that at least 2.5 million peasants (in addition to the above industrial workers) were deported, though the true number is believed to be much greater. The Great Purges and Politics in the Soviet Union (1930-1939). Throughout the period of collectivization and the Five-Year Plan, the Soviet government became increasingly tyrannical. Stalin, was extremely paranoid, began to turn on important members of the party he had once called supporters. In 1934, the last person who might have rivaled Stalin, Sergei Kirov, was shot in his office, most likely on Stalin's orders. Using the murder as a pretext, he began to engage in ruthless purges of the party membership. Ironically, most of the purged members were original members of the party and colleagues of Lenin, known as the Old Bolsheviks. Through a series of show trials, the defendants were sentenced to death and to forced labor in the Gulag. Often, after using torture to extract signed confessions and agreeing on lenient sentences for a confession of false charges in the court, Stalin would turn on his word and have the defendants executed. Zinoviev and Kamenev, Stalin's old allies, both met this fate. Through 1936-1937, a period known as the Great Terror, Stalin supposedly personally signed 40,000 death warrants. Stalin's dictatorship held incredible control over the general populace of the nation. Intense propaganda campaigns tried to indoctrinate the society with Communist thought. Stalin wanted to replace the national identities, such as the Russians, Ukrainians, or Belarussians, with the idea of a purely "Soviet" citizen. He also stipulated that all ethnic groups be treated equally. Under the Tsars, the Russians had been given preference. Now the heavy hand of Stalin was given equally to all nations. However, this did not prevent him from forcing the speakers of every language in the USSR to convert to the Cyrillic alphabet. Religion came under intense pressure as well, as atheism was the official policy of the state. Priests were rounded up and shipped to Gulag or executed. By the end of the terror, less than 1,000 churches remained out of at least 20,000. The NKVD, the Soviet secret police, hunted down citizens suspected of "counter-revolutionary" or "subversive" crimes. During the Great Terror, as many as 1 million people (the NKVD's own records admit 0.681 million) were executed for simply "opposing" Stalin's ideas and plans. False confessions were routinely extracted through torture and intimidation. The penalty for countless others (numbering by the most conservative estimates in the millions) was the Gulag. Fear was the order of the day in the Soviet Union. In addition, the Soviet state cultivated an extreme cult of personality around Stalin. Pictures of the dictator appeared at every street corner and in every building, including people's homes. School children ended the pledge of allegiance at the beginning of each day by saying "...and thank Comrade Stalin for this happy life". To be fair, social conditions did improve under Stalin. Unemployment fell to practically zero, and large bounds in the public health were introduced. However, there was no freedom whatsoever in Soviet society. Francisco Franco's Fascist Spain and the Spanish Civil War (1936-1945). Francisco Franco rose to power after the brutal turmoil of the Spanish Civil War. The conflict was between the residing leftist republican parties and the nationalist movement led by Franco. Some view the Spanish Civil War as a "dress rehearsal for World War II". The Struggle of dictatorship and democracy is evident in this conflict. 

The First World War left Europe in ruins. By the Treaty of Versailles, Germany, taking the blame for the war, was forced to pay massive reparations to the victorious Allies and lost lands to France and a restored Poland. The Austro-Hungarian Empire collapsed, with Czechoslovakia seceding, and other territories going to Poland, Serbia (now Yugoslavia), Italy, and Romania. These changes were made permanent by the Treaty of Saint-Germain, which also gave Yugoslavia the Austrian Adriatic Fleet. Bulgaria ceded small strips of territory to Romania, Yugoslavia, and Greece. Most devastating, though, was a massive influenza outbreak that started on the front lines and spread throughout the world, carried by soldiers returning from the war. This outbreak killed more people in more countries than the war itself. 

Causes of World War II. France, Great Britain, and the U.S. had attained their wartime objectives in 1919. They had reduced Germany to a military cipher and had reorganized Europe and the world as they saw fit. The French and the British frequently disagreed on policy in the postwar period, however, and were unsure of their ability to defend the peace settlement. Disillusionment with war led to the practice of appeasement, or giving into an aggressor's demands to keep the peace. The U.S., disillusioned by the Europeans' failure to pay their war debts, retreated into isolationism. The Treaty of Versailles left many countries dissatisfied. Adverse conditions, such as reparations and unemployed veterans from World War I led to the circulation of new, radical ideas and solutions, such as fascism in Italy. This Fascist party, as Mussolini called it, later became a model for Hitler in Germany. The Failure of Peace Efforts. During the 1920s, attempts were made to achieve a stable peace. The first was President Woodrow Wilson's idea to establish the League of Nations (1920) as a forum in which nations could settle their disputes. The League's powers were limited to persuasion and various levels of moral and economic sanctions that the members were free to carry out as they saw fit. The United States never joined the League and Germany and the USSR were also never members. At the Washington Conference of 1921-2, the principal naval powers agreed to limit their navies according to a fixed ratio. The Locarno Conference (1925) produced a treaty guarantee of the German-French boundary and an arbitration agreement between Germany and Poland. In the Kellogg-Briande Pact (1928), 63 countries including all the Great Powers except the USSR, renounced war as an instrument of national policy and pledged to resolve all disputes among them "by pacific means." The signatories had agreed beforehand to exempt wars of "self-defense." The Rise of Fascism. One of the victors' stated aims in World War I had been "to make the world safe for democracy," and postwar Germany adopted a democratic constitution, as did most of the other states restored or created after the war. In the 1920s, however, the wave of the future appeared to be a form of nationalistic, militaristic totalitarianism known by its Italian name, fascism. It promised to minister to peoples' wants more effectively than democracy and presented itself as the one sure defense against communism. Benito Mussolini established the first Fascist, European dictatorship during the inter war period in Italy in 1922. Formation of the Axis Coalition. Adolf Hitler, the Leader of the German National Socialist (Nazi) party, preached a racist brand of fascism. Hitler promised to overturn the Versailles Treaty and secure additional "Lebensraum" ("living space") for the German people, who he contended deserve more as members of a superior race. In the early 1930s, the Great Depression hit Germany. The moderate parties could not agree on what to do about it, and large numbers of voters turned to the Nazis and Communists. In 1933 Hitler became the German Chancellor, and in a series of subsequent moves established himself as dictator. Japan did not formally adopt fascism, but the armed forces' powerful position in government enabled them to impose a similar type of totalitarianism. As dismantlers of the world status quo, the Japanese were well ahead of Hitler. They used a minor clash with Chinese troops near Mukden, also known as the Mukden or Manchurian crisis, in 1931 as a pretext for taking over all of Manchuria, where they proclaimed the puppet state of Manchukuo in 1932. In 1937-8 they occupied the main Chinese ports. Having denounced the disarmament clauses of the Versailles Treaty, created a new air force, and reintroduced conscription, Hitler tried out his new weapons on the side of right-wing military rebels in the Spanish civil war (1936-9). This venture brought him into collaboration with Mussolini who was also supporting the Spanish revolt after having seized (1935-6) Ethiopia in a small war. Treaties between Germany, Italy, and Japan in 1936-7 brought into being the Rome-Berlin-Tokyo Axis. For example, Japan and Germany signed the Anti-Comintern pact in 1936 and then Italy joined in 1937. This pact denounced communism and it showed their unity in the matter. The Axis thereafter became the collective term for those countries and their allies. German Aggression in Europe. Hitler launched his own expansionist drive with the annexation of Austria in March 1938. The way was clear: Mussolini supported him; and the British and French, overawed by German rearmament, accepted Hitler's claim that the status of Austria was an internal German affair. The U.S. had impaired its ability to act against aggression by passing a neutrality law that prohibited material assistance to all parties in foreign conflicts. In September 1938 Hitler threatened war to annex the western border area of Czechoslovakia, the Sudetenland and its 3.5. million ethnic Germans. The British Prime Minister Neville Chamberlain initiated talks that culminated at the end of the month in the Munich Pact, by which the Czechs, on British and French urging, relinquished the Sudetenland in return for Hitler's promise not to take any more Czech territory. Chamberlain believed he had achieved "peace for our time," but the word Munich soon implied abject and futile appeasement. Less than six months later, in March 1939, Hitler seized the remainder of Czechoslovakia. Alarmed by this new aggression and by Hitler's threats against Poland, the British government pledged to aid that country if Germany threatened its independence. A popular joke ran at the time: "A guarantee a day keeps Hitler away". France already had a mutual defense treaty with Poland. The turn away from appeasement brought the Soviet Union to the fore. Joseph Stalin, the Soviet dictator, had offered military help to Czechoslovakia during the 1938 crisis, but had been ignored by all the parties to the Munich Agreement. Now that war threatened, he was courted by both sides, but Hitler made the more attractive offer. Allied with Britain and France, the Soviet Union might well have had to fight, but all Germany asked for was its neutrality. In Moscow, on the night of August 23, 1939, the Nazi-Soviet Pact was signed. In the part published the next day, Germany and the Soviet Union agreed not to go to war against each other. A secret protocol gave Stalin a free hand in Finland, Estonia, Latvia, eastern Poland, and eastern Romania. The Worldwide Great Depression. The costs of carrying out World War I, as well as the costs to rebuild Western Europe after years of fighting, resulted in enormous debts on the part of the Western European powers to the United States. The enormous reparations put on Germany in the Treaty of Versailles also increased the debts. Coupled with ineffective governments in many of these European States (notably the Weinmar Republic, pre-Mussolini Italy and Socialist France) led to slow reconstruction and poor economic growth. With the crash of the New York Stock Market on 29 October, 1929, the United States recalled all foreign loans in the following days. Unable to repay these loans, the economies of the West collapsed, beginning the Great Depression. War in Europe. The German invasion of Poland on September 1, 1939 is regarded as the start of World War II. Due to official League of Nations law, France and the UK were obligated to intervene, but in reality did very little in a period known as the phony war. In Poland, however, fighting had begun. Germany was extremely successful in this war due to their new strategy known as blitzkrieg (meaning "lightning war"). This strategy involved the use of tanks, a relatively new technology at the time, to quickly overrun a country before they could set up their defenses properly. This strategy worked amazingly for Germany, with the German troops arriving near the Polish capitol of Warsaw in just 7 days. Germany and the USSR had made a nonaggression pact in the days leading up to the invasion, which was essentially an agreement to split eastern Europe between the two countries. The agreement was initially a shock to the world, as the two countries were governed by ideologies that were about as far apart from each other as possible. However, this fragile agreement did not last. Before that, though, the USSR used this opportunity to annex the Baltic states, parts of Romania, and tried to take land Finland in a disastrous campaign for the Russians known as the winter war. Germany then proceeded to take out Denmark and Norway in order to secure positions in the north sea. Then Germany shifted their attention to France. French strategy relied on using the Maginot line, a massive line of forts stretching the entirety of their border with Germany to protect their border with relatively few troops while positioning their best forces along the border with Belgium, as Germany had attacked through Belgium in the first world war. However, they left the Ardennes forest undefended, as they believed that the dense forest would provide enough of a natural barrier. Knowing this, Germany launched a daring attack that sent 50 tank divisions through the Ardennes. The Germans broke through, and managed to cut off the allied forces. The French forces there were wiped out, but they were able to defend long enough to let British forces escape through Dunkirk. But for France, all was lost. The bulk of their army had been wiped out, and within weeks France had to surrender. This was far from the end of French resistance, though, as aside from civilian resistance efforts, and forces fighting under the banner of "Free France" played a major role in the north African front in the coming years. In Europe, however, Germany had set up a puppet state called Vichy France in the south, while the north was fully annexed for defense purposes. Germany tried to make the odds seem as far in their favor as they could in the hopes that Britain would try to make peace, but the new Prime Minister Winston Churchill refused to give up. In order to take over the British Isles, Germany would first have to gain control of the English channel in both the air and sea to be able to land troops on the island. After intense air fighting over England that resulted in the bombing of major English cities, Britain stopped any chance of a naval invasion. The War in the Pacific. Mukden Incident and the Invasion of Manchuria (1931). After winning the Russo-Japanese War in 1905, Japan quickly became the dominant power in its region. Russia recognized Korea as a Japanese sphere of influence and removed all of its forces from there and Manchuria, the sparsely populated northeastern region of China. In 1910, Japan annexed Korea as its own with little protest or resistance. Still, Japan was a quickly growing country, both population-wise and economically. It founded the South Manchuria Railway company in Manchuria in 1906, and with that company was able to gain government-like control of the area. By 1931, the Depression had struck a blow to Japan. The government did little to help Japan's economy, and in the eyes of its citizens, was weak and powerless. Instead, the public favored the Japanese army, and soon the civilian government had lost control of its military. To the army, Manchuria seemed like an obvious solution to many of Japan's problems. Manchuria was vast and thinly populated, and would serve as excellent elbow room for an already overcrowded Japan. It was also thought that Manchuria was rich in forests, natural resources, and fertile land. The fact that the Japanese believed themselves to be far superior to the Chinese only moved Japan towards conflict faster. Additionally, the warlord of Manchuria went against Japanese expectations and declared his allegiance to a growing Chinese military movement. So, in 1931, the army staged an explosion at a section of railway near Mukden, a city in Manchuria, as a pretext to invade and annex China. Japan met little resistance, although it did not have support of its own government, and Manchuria was completely occupied by the end of the year. Japan subsequently set up the puppet state of Manchukuo to oversee the newly acquired region. The League of Nations vehemently protested Japan's aggression, but Japan then withdrew from it. Japan invades China (1937). The 1920s saw a weak and politically chaotic China. Warlords of the many provinces of China constantly feuded, and the central government was weak and decentralized, unable to do anything to stop conflict. In 1927 Chiang Kai-Shek gained control of the Kuomintang (the Chinese government) and its National Revolution Army. Chiang led an expedition to defeat southern and central Chinese warlords and gain the allegiance of northern warlords. He was successful, and he soon focused on what he perceived to be a greater threat than Japan, which was communism. But in 1937, the deposed warlord general of Manchuria kidnapped Chiang and refused to release him until he at least temporarily united with the communists against the Japanese threat. The Japanese army responded by staging the Battle of Lugou Bridge, which was supposed to provoke open war between China and Japan. It worked and the Sino-Japanese War began. The beginning of the conflict was marked by the Chinese strategy of giving up land in order to stall the Japanese. It is important to note that the Japanese was not to completely take over China; rather, the Japanese wanted to set up puppet governments in key regions that would protect and advance Japanese interests. The fall of Nanjing in the early stages of this conflict saw the beginning of Japanese war atrocities. 100,000-300,000 were killed in the six weeks after Nanjing was captured. Other war crimes committed included widespread rape, arson, and looting. Anti-Comintern Pact and Tripartite Pact. These were pacts between Germany, Italy, and Japan. The Anti-Comintern pact had been a pact that denounced communism and it was initially signed by Japan and Germany. However, later, as German and Italian relations improved, Italy also signed and this was made stronger later by the Rome-Berlin-Tokyo Axis in 1938. The Tripartite Pact also strengthened the alliance and it was basically a confirmation of the Rome-Berlin-Toyko Axis. Pearl Harbor and Simultaneous Invasions (early December 1941). On December 7, 1941, Japanese warplanes commanded by Vice Admiral Chuichi Nagumo carried out a surprise air raid on Pearl Harbor, Hawaii, the largest U.S. naval base in the Pacific. The Japanese forces met little resistance and devastated the harbor. This attack resulted in 8 battleships either sunk or damaged, 3 light cruisers and 3 destroyers sunk as well as damage to some auxiliaries and 343 aircraft either damaged or destroyed. 2408 Americans were killed including 68 civilians; 1178 were wounded. Japan lost only 29 aircraft and their crews and five midget submarines. However, the attack failed to strike targets that could have been crippling losses to the US Pacific Fleet such as the aircraft carriers which were out at sea at the time of the attack or the base's ship fuel storage and repair facilities. The survival of these assets have led many to consider this attack a catastrophic long term strategic blunder for Japan. The following day, the United States declared war on Japan. Simultaneously to the attack on Pearl Harbor, Japan also attacked U.S. air bases in the Philippines. Immediately following these attacks, Japan invaded the Philippines and also the British Colonies of Hong Kong, Malaya, Borneo and Burma with the intention of seizing the oilfields of the Dutch East Indies. Following the Japanese attack on Pearl Harbor, Germany declared war on the United States on December 11, 1941, even though it was not obliged to do so under the Tripartite Pact of 1940. Hitler made the declaration in the hopes that Japan would support him by attacking the Soviet Union. Japan did not oblige him, and this diplomatic move proved a catastrophic blunder which gave President Franklin D. Roosevelt the pretext needed for the United States joining the fight in Europe with full commitment and with no meaningful opposition from Congress. Some historians mark this moment as another major turning point of the war with Hitler provoking a grand alliance of powerful nations, most prominently the UK, the USA and the USSR, who could wage powerful offensives on both East and West simultaneously. Allied Defeats in the Pacific and Asia (late December 1941-1942). Simultaneous with the dawn raid on Pearl Harbor, the Japanese carried out an invasion of Malaya, landing troops at Kota Bharu on the east coast, supported by land based aircraft from bases in Vietnam and Taiwan. The British attempted to oppose the landings by dispatching Force Z, comprising the battleship HMS Prince of Wales and the battlecruiser HMS Repulse, with their escorting destroyers, from the naval base in Singapore, but this force was intercepted and destroyed by bombers before even reaching their objective. In a series of swift maneuvers down the Malay peninsula, thought by the British to be "impassable" to an invading force landing so far north, the Japanese advanced down to the Johor Straits at the southernmost tip of the peninsula by January 1942. The Japanese were even using tanks, which the British had thought would not be able to penetrate the jungles but they were wrong. During a short two week campaign the Japanese crossed the Straits of Johor by amphibious assault and conducted a series of sharp battles, notably the battle of Kent Ridge when the Royal Malay Regiment put up a brave but futile effort to stem the tide. Singapore fell on February 15, 1942 and with its fall, Japan was now able to control the sea approaches from the Indian Ocean through the Malacca Straits. The natural resources of the Malay peninsula, in particular rubber plantations and tin mines, were now in the hands of the Japanese. Other Allied possessions, especially in the oil rich East Indies (Indonesia) were also swiftly captured, and all organised resistance effectively ceased, with attention now shifting to events closer to Midway, the Solomon Islands, the Bismark Sea and New Guinea. Allies Regroup and the Battle of Midway (1942). Following the attack on Pearl Harbour, the US military sought to strike back at Japan, and a plan was formulated to bomb Tokyo. As Tokyo could not be reached by land based bombers, it was decided to use an aircraft carrier to launch the attack close to Japanese waters. The Doolittle Raid was carried out by James H. (Jimmy) Doolittle and his squadron of B-25 medium bombers, launched from the USS Hornet. The raid achieved little strategically, but was a tremendous morale booster in the dark days of 1942. It also led to the decision by the Japanese military to attack the only logical base of the attackers, the tiny atoll of Midway. A powerful force of warships, with four large fleet carriers at its core (Akagi, Kaga, Hiryu and Soryu) attacked Midway. The US navy, with the aid of intercepted and decoded Japanese signals, were ready and launched a counter attack with the carriers USS Enterprise and USS Yorktown, destroying all four of the Japanese fleet carriers. This was a devastating blow to the Japanese and is considered the turning point of the Pacific War. The Japanese had largely roamed the Pacific Ocean, the South China Sea, the Malacca Straits and the Indian Ocean with impunity, launching raids from these same four carriers on Allied bases in these areas including Darwin, Colombo and along the Indian east coast. With the loss of these carriers and more importantly their cadre of irreplaceable hard core highly trained naval aviators, the Japanese could no longer maintain an effective offensive and became largely defensive from then on. Island Hopping (1943- Late 1944). Island hopping was a campaign of capturing key islands in the Pacific that were used as prerequisites, or stepping stones, to the next island with the eventual destination being Japan, rather than trying to capture every island under Japanese control. Allied forces often assaulted weaker islands first, while starving out the Japanese strongholds before attacking them. The Atomic Bomb (August 1945). On August 6, 1945, a lone B-29 bomber, named the Enola Gay, appeared over the skies of Hiroshima. Air raid sirens went off around the city and people ran for their shelters. However, minutes later, the all-clear symbol was given. Although it had been a seemingly harmless run, the B-29 had, in fact, dropped a single bomb (this bomb was called "Little Boy"). This bomb detonated about 1,900 feet over Hiroshima and leveled much of the city within a few thousandths of a second. Tens of thousands were killed immediately and many more would eventually die from the radiation poisoning. However, Japan did not surrender to the United States, so three days later, on August 9, 1945, a B-29 named Boxcar dropped an atom bomb on the city of Nagasaki (this bomb was called "Fat Man"). Although the bomb was actually more powerful than the Hiroshima bomb, the foggy weather conditions and the hilly terrain of Nagasaki somewhat shielded a portion of the city from the worst effects. This led to an immediate ceasefire with Japan, and surrender a month later. 

The Second World War saw the most far-reaching transformation of world politics to date. The destructive technologies introduced during the war – foremost, the atomic bomb – made it very unlikely that a land-based conflict of similar scale and duration among the major nations could ever happen again, because of the potential for total destruction of all combatants. No advanced industrial nation has been invaded since 1945, and all wars since that time have been guerrilla conflicts in less-developed countries, conflicts involving less-developed countries with more advanced ones, or some combination of these two scenarios. From an economic standpoint, the war and its aftermath consumed much of the real and potential industrial production of the world over the period 1940–1960 (with the exception that the United States, its homeland untouched, was able to expand both its defense industries and its civilian economy very rapidly after 1945). Europe and Japan lay in ruins and would spend 15–20 years rebuilding the basis for economic life, with much assistance from the U.S. The Soviet Union and China, though victorious in the war, were also ravaged. Splitting of the world. Europe was split into two main camps by the "Iron Curtain", which divided Germany in half and separated Austria from Czechoslovakia and Hungary, and Italy from Yugoslavia. The Soviet Union absorbed eastern Poland, and "reassigned" large areas of German territory to Polish rule by way of compensation. Moscow intervened directly to install Communist parties in power in Poland, eastern Germany, Hungary and Czechoslovakia. Finland was able to keep its independence, but did not regain the lands it lost to the Soviets in the 1940 Winter War. Yugoslavia under Marshal Tito, already Communist, did not submit to direct influence from Moscow, choosing a more independent path and greatly angering Stalin. Elsewhere in the Balkans, Bulgaria, Romania and Albania also were brought into the Soviet bloc. To the west were the democracies allied to the USA: the UK, France, Italy, the Netherlands, Belgium, Norway, and West Germany. Washington became quite concerned, however, that local Communist parties might gain power in France, Italy and Greece in the late 1940s, given the battered state of the postwar European economy and the proximity of the Soviet Union. French leader Charles DeGaulle received strong backing from the U.S., and anti-Communist parties in Italy were heavily financed. In the end neither nation left the Western sphere of influence. In Asia, China, Vietnam, Korea and Mongolia were communist countries. Japan, first under occupation by the USA after WWII, had to reform its system – moving away from from militarism and expansionism and into democratic reforms. Military alliances were formed on both sides: First NATO, on April 4, 1949 in the USA geopolitical sphere and as a direct response due to the Cold War the Soviet Union (USSR), on May 14, 1955 created the Organization of Warsaw treaty, better known as Warsaw Pact. In an attempt to set up a body for possible dialog between Western, Eastern, and Developing countries and establish international law, the Organization of United Nations was established in 1945 by USA initiative. The failing of the League of Nations, also initiated by the USA was one of several precursor stages for WWII. Special consideration must be given to the geopolitical interests behind these organizations, especially who they favor. Even if ultimately, they are a stabilizing force. Other USA international initiatives include: the International Monetary Fund (IMF), the World Bank and more recently the World Trade Organization (WTO), that has its roots in the discussions to form the International Trade Organization (ITO) as an evolution from the General Agreement on Tariffs and Trade (GATT) of 1947. 

Grammar - Object Pronouns. Direct Object Pronouns. While the subject of a sentence "initiates" an action (the verb), the direct object is the one that is "affected" by the action. A direct object pronoun is used to refer to the direct object of a previous sentence: The following table shows the six types of direct object pronouns: In spanish tú is used for informal situations, and usted must be used when a formal treatment is needed. Note: In Spain, "le" and "les" are used as the masculine direct object pronoun only when referring to people. If the antecedent of a direct object is masculine but non-human, "lo" or "los" are used instead. In most other Spanish speaking places, "lo" and "los" are used instead of "le" and "les". Indirect Object Pronouns. An indirect object is an object that would be asked for with "To whom...?" or "From whom...?". It is called "indirect" because it occurs usually together with a direct object which is affected directly by the action: The apple "is given" by the woman (direct). The boy gets the "given apple" (indirect - depends on the apple being given). Here is a table with all of the Spanish indirect object pronouns: Position Of Object Pronouns (Double Object Pronouns). So far we have only seen sentences with one object pronoun. If there is both a direct and an indirect object pronoun, the indirect pronoun usually comes first: Also, when both object pronouns are in the third person (either singular or plural), the indirect pronoun changes from le/les to se: In sentences that contain an infinitive or a participle, the object pronoun may be either placed before the conjugated verb or it maybe attached to the infinitive/participle: It is possible to have the two rules above working at the same time: A combination of direct and indirect pronouns that is attached to an infinitive/participle: Exercise:Object Pronouns Vocabulario (Vocabulary) - La comida (Food).  In Spain and some other countries, "comida" is the midday meal.  In other countries, for example Chile, "comida" is the last meal in the day. Instead of saying "desayuno, comida y cena" (Spain) or "desayuno, almuerzo y comida" (Chile, Colombia), it's safer to say "desayuno, almuerzo y cena". The word "comida" has several meanings Note that due to the pervasive influence of English, in many supermarkets there is a section called "Vegetales" instead of "Verduras". They mistranslate vegetable, forgetting that this is not the same as English vegetal (relating to plants). 

These books describe and explain Perl, a high-level, general-purpose, interpreted, dynamic programming language. Perl is noted for its idiomatically rich syntax, its extensive use of regular expressions and the large module archive CPAN. 





As you have learned in the previous sections, the linear waves you studied in the previous section tend to disperse. However, there are nonlinear waves described by nonlinear partial differential equations (PDE) which admits solutions with nondispersing wave packets. Here are two primary examples: Sine-Gordon. This is described by the PDE formula_1. The equation of a wave travelling along the positive x direction is given by y=0.25 x 0.001sin(500t-0.025x)determine the angular frequence 

What is Differentiation? Differentiation is a process of finding a function that outputs the rate of change of one variable with respect to another variable. Informally, we may suppose that we're tracking the position of a car on a two-lane road with no passing lanes. Assuming the car never pulls off the road, we can abstractly study the car's position by assigning it a variable, formula_1 . Since the car's position changes as the time changes, we say that formula_1 is dependent on time, or formula_3 . This tells where the car is at each specific time. Differentiation gives us a function formula_4 which represents the car's speed, that is the rate of change of its position with respect to time. Equivalently, differentiation gives us the slope at any point of the graph of a non-linear function. For a linear function, of form formula_5 , formula_6 is the slope. For non-linear functions, such as formula_7 , the slope can depend on formula_1 ; differentiation gives us a function which represents this slope. The Definition of Slope. Historically, the primary motivation for the study of differentiation was the tangent line problem: for a given curve, find the slope of the straight line that is tangent to the curve at a given point. The word "tangent" comes from the Latin word "tangens", which means touching. Thus, to solve the tangent line problem, we need to find the slope of a line that is "touching" a given curve at a given point, or, in modern language, that has the same slope. But what exactly do we mean by "slope" for a curve? The solution is obvious in some cases: for example, a line formula_9 is its own tangent; the slope at any point is formula_10 . For the parabola formula_11 , the slope at the point formula_12 is formula_13 ; the tangent line is horizontal. But how can you find the slope of, say, formula_14 at formula_15 ? This is in general a nontrivial question, but first we will deal carefully with the slope of lines. The Slope of a Line. The slope of a line, also called the gradient of the line, is a measure of its inclination. A line that is horizontal has slope 0, a line from the bottom left to the top right has a positive slope and a line from the top left to the bottom right has a negative slope. The slope can be defined in two (equivalent) ways. The first way is to express it as how much the line climbs for a given motion horizontally. We denote a change in a quantity using the symbol formula_16 (pronounced "delta"). Thus, a change in formula_1 is written as formula_18 . We can therefore write this definition of slope as: An example may make this definition clearer. If we have two points on a line, formula_20 and formula_21 , the change in formula_1 from formula_23 to formula_24 is given by: Likewise, the change in formula_26 from formula_23 to formula_24 is given by: This leads to the very important result below. Alternatively, we can define slope trigonometrically , using the tangent function: where formula_31 is the angle from the rightward-pointing horizontal to the line, measured counter-clockwise. If you recall that the tangent of an angle is the ratio of the y-coordinate to the x-coordinate on the unit circle, you should be able to spot the equivalence here. Of a graph of a function. The graphs of most functions we are interested in are not straight lines (although they can be), but rather curves. We cannot define the slope of a curve in the same way as we can for a line. In order for us to understand how to find the slope of a curve at a point, we will first have to cover the idea of tangency. Intuitively, a tangent is a line which "just" touches a curve at a point, such that the angle between them at that point is 0. Consider the following four curves and lines: A secant is a line drawn through two points on a curve. We can construct a definition of a tangent as the limit of a secant of the curve taken as the separation between the points tends to zero. Consider the diagram below. As the distance formula_32 tends to 0, the secant line becomes the tangent at the point formula_33 . The two points we draw our line through are: and As a secant line is simply a line and we know two points on it, we can find its slope, formula_36 , using the formula from before: (We will refer to the slope as formula_36 because it may, and generally will, depend on formula_32 .) Substituting in the points on the line, This simplifies to This expression is called the difference quotient. Note that formula_32 can be positive or negative — it is perfectly valid to take a secant through any two points on the curve — but cannot be formula_13 . The definition of the tangent line we gave was not rigorous, since we've only defined limits of "numbers" — or, more precisely, of functions that output numbers — not of "lines". But we "can" define the "slope" of the tangent line at a point rigorously, by taking the limit of the slopes of the secant lines from the last paragraph. Having done so, we can "then" define the tangent line as well. Note that we cannot simply set formula_32 to 0 as this would imply division of 0 by 0 which would yield an undefined result. Instead we must find the limit of the above expression as formula_32 tends to 0: This last equation is just the point-slope form for the line through formula_46 with slope formula_10. Exercises. Solutions The Rate of Change of a Function at a Point. Consider the formula for average velocity in the formula_1 direction, formula_49 , where formula_18 is the change in formula_1 over the time interval formula_52 . This formula gives the average velocity over a period of time, but suppose we want to define the instantaneous velocity. To this end we look at the change in position as the change in time approaches 0. Mathematically this is written as: formula_53 , which we abbreviate by the symbol formula_4 . (The idea of this notation is that the letter formula_55 denotes change.) Compare the symbol formula_55 with formula_16 . The idea is that both indicate a difference between two numbers, but formula_16 denotes a finite difference while formula_55 denotes an infinitesimal difference. Please note that the symbols formula_60 and formula_61 have no rigorous meaning on their own, since formula_62 , and we can't divide by 0. The Definition of the Derivative. You may have noticed that the two operations we've discussed — computing the slope of the tangent to the graph of a function and computing the instantaneous rate of change of the function — involved exactly the same limit. That is, the slope of the tangent to the graph of formula_67 is formula_68 . Of course, formula_68 can, and generally will, depend on formula_1 , so we should really think of it as a "function" of formula_1 . We call this process (of computing formula_68) differentiation. Differentiation results in another function whose value for any value formula_1 is the slope of the original function at formula_1 . This function is known as the derivative of the original function. Since lots of different sorts of people use derivatives, there are lots of different mathematical notations for them. Here are some: Most of the time the brackets are not needed, but are useful for clarity if we are dealing with something like formula_83 , where we want to differentiate the product of two functions, formula_84 and formula_85 . The first notation has the advantage that it makes clear that the derivative is a function. That is, if we want to talk about the derivative of formula_76 at formula_87 , we can just write formula_88 . In any event, here is the formal definition: Examples. Example 1 The derivative of formula_89 is no matter what formula_1 is. This is consistent with the definition of the derivative as the slope of a function. Example 2 What is the slope of the graph of formula_92 at formula_93 ? We can do it "the hard (and imprecise) way", "without" using differentiation, as follows, using a calculator and using small differences below and above the given point: When formula_94 , formula_95 . When formula_96 , formula_97 . Then the difference between the two values of formula_1 is formula_99 . Then the difference between the two values of formula_26 is formula_101 . Thus, the slope formula_102 at the point of the graph at which formula_103 . But, to solve the problem precisely, we compute We were lucky this time; the approximation we got above turned out to be exactly right. But this won't always be so, and, anyway, this way we didn't need a calculator. In general, the derivative of formula_7 is Example 3 If formula_105 (the absolute value function) then formula_106 , which can also be stated as Finding this derivative is a bit complicated, so we won't prove it at this point. Here, formula_76 is not smooth (though it is continuous) at formula_109 and so the limits formula_110 and formula_111 (the limits as 0 is approached from the right and left respectively) are not equal. From the definition, formula_112 , which does not exist. Thus, formula_113 is undefined, and so formula_75 has a discontinuity at 0. This sort of point of non-differentiability is called a cusp. Functions may also not be differentiable because they go to infinity at a point, or oscillate infinitely frequently. Understanding the derivative notation. The derivative notation is special and unique in mathematics. The most common notation for derivatives you'll run into when first starting out with differentiating is the Leibniz notation, expressed as formula_68 . You may think of this as "rate of change in formula_26 with respect to formula_1" . You may also think of it as "infinitesimal value of formula_26 divided by infinitesimal value of formula_1" . Either way is a good way of thinking, although you should remember that the precise definition is the one we gave above. Often, in an equation, you will see just formula_120 , which literally means "derivative with respect to x". This means we should take the derivative of whatever is written to the right; that is, formula_121 means formula_68 where formula_123 . As you advance through your studies, you will see that we sometimes pretend that formula_124 and formula_60 are separate entities that can be multiplied and divided, by writing things like formula_126 . Eventually you will see derivatives such as formula_127 , which just means that the input variable of our function is called formula_26 and our output variable is called formula_1 ; sometimes, we will write formula_130 , to mean the derivative with respect to formula_26 of whatever is written on the right. In general, the variables could be anything, say formula_132 . All of the following are equivalent for expressing the derivative of formula_11 Exercises. Solutions Differentiation Rules. The process of differentiation is tedious for complicated functions. Therefore, rules for differentiating general functions have been developed, and can be proved with a little effort. Once sufficient rules have been proved, it will be fairly easy to differentiate a wide variety of functions. Some of the simplest rules involve the derivative of linear functions. Derivative of a constant function. For any fixed real number formula_139 , formula_140 Intuition. The graph of the function formula_141 is a horizontal line, which has a constant slope of 0. Therefore, it should be expected that the derivative of this function is zero, regardless of the values of formula_1 and formula_139 . Proof. The definition of a derivative is Let formula_141 for all formula_1 . (That is, formula_84 is a constant function.) Then formula_148 . Therefore Let formula_150 . To prove that formula_151 , we need to find a positive formula_152 such that, for any given positive formula_153 , formula_154 whenever formula_155 . But formula_156 , so formula_154 for any choice of formula_152 . Examples. Note that, in the second example, formula_161 is just a constant. Derivative of a linear function. For any fixed real numbers formula_10 and formula_139 , formula_164 The special case formula_165 shows the advantage of the formula_120 notation—rules are intuitive by basic algebra, though this does not constitute a proof, and can lead to misconceptions to what exactly formula_60 and formula_124 actually are. Intuition. The graph of formula_9 is a line with constant slope formula_10. Proof. If formula_171 , then formula_172. So, Constant multiple and addition rules. Since we already know the rules for some very basic functions, we would like to be able to take the derivative of more complex functions by breaking them up into simpler functions. Two tools that let us do this are the constant multiple rule and the addition rule. The Constant Rule. For any fixed real number formula_139 , formula_174 The reason, of course, is that one can factor formula_139 out of the numerator, and then of the entire limit, in the definition. The details are left as an exercise. Example We already know that Suppose we want to find the derivative of formula_177 Another simple rule for breaking up functions is the addition rule. The Addition and Subtraction Rules. formula_178 Proof From the definition: By definition then, this last term is formula_179 Example What is the derivative of formula_180 ? The fact that both of these rules work is extremely significant mathematically because it means that differentiation is linear. You can take an equation, break it up into terms, figure out the derivative individually and build the answer back up, and nothing odd will happen. We now need only one more piece of information before we can take the derivatives of any polynomial. The Power Rule. formula_181 This has been proved in an example in Derivatives of Exponential and Logarithm Functions where it can be best understood. For example, in the case of formula_182 the derivative is formula_183 as was established earlier. A special case of this rule is that formula_184 . Since polynomials are sums of monomials, using this rule and the addition rule lets you differentiate any polynomial. A relatively simple proof for this can be derived from the binomial expansion theorem. This rule also applies to fractional and negative powers. Therefore Derivatives of polynomials. With these rules in hand, you can now find the derivative of any polynomial you come across. Rather than write the general formula, let's go step by step through the process. The first thing we can do is to use the addition rule to split the equation up into terms: We can immediately use the linear and constant rules to get rid of some terms: Now you may use the constant multiplier rule to move the constants outside the derivatives: Then use the power rule to work with the individual monomials: And then do some algebra to get the final answer: These are not the only differentiation rules. There are other, more advanced, differentiation rules, which will be described in a later chapter. Exercises. Solutions 

Fourier Transform. So far, you've learned how to superimpose a finite number of sinusoidal waves. However, a wave in general can't be expressed as the sum of a finite number of sines and cosines. Fortunately, we have a theorem called Fourier's theorem which basically states that under certain technical assumptions, any function, f(x) is equal to an integral over sines and cosines. In other words, Now, if we're given the wave function when t=0, φ(x,0) and the velocity of each sine wave as a function of its wave number, v(k), then we can compute φ(x,t) for any t by taking the inverse Fourier transform of φ(x,0) conducting a phase shift, and then taking the Fourier transform. Fortunately, the inverse Fourier transform is very similar to the Fourier transform itself. This tells us that, since waves which are very spread out, like the sine wave, have a narrow range of wave numbers, wave functions whose wave numbers are very spread out will only be significant at a narrow range of positions. 

Writing General Chemistry Reactions. In organic chemistry, a reaction may be written precisely as it is for general chemistry if only a basic amount of information is needed. For example, when a haloalkane is turned into an alkene, the reaction may be written: codice_1 Unfortunately, this method of notation does not tell anyone very much about the reaction, and it takes expertise to know exactly what is going on. A new student to organic chemistry probably would not notice that the product molecule contains one site of unsaturation due to a double bond between carbon atoms number one and number two. Because it is so general, this notation is good for general chemistry, but organic chemistry requires more precision. For most students, common practices in writing organic reactions will be different than used in general chemistry. Differences in Organic Chemistry Notation. Organic chemistry reactions are often not written as balanced equations. This is because many organic chemists - who are just as lazy as anyone else - tend to be more interested in the "organic product" of a reaction than in anything else going on in the reaction. Side products are often ignored, and just as often catalysts and solution notation may be highly abbreviated or left out altogether. As you gain familiarity with organic chemistry you will come to understand just what may be abbreviated or left out, but in the beginning this can be a source of frustration. Another difference is that modified Lewis drawings of molecules are often used instead of molecular formulas. This makes sense due to the fact that organic molecules are often rather large in size and complicated in structure, so that they can be more easily understood in the form of a drawing as opposed to a word-formula. A two-dimensional drawing reveals some of the three-dimensional shape of the molecule, but when necessary even three-dimensional drawings are used to depict reactants and products. &lt;br&gt; &lt;br&gt; Working with the above drawing of a molecule may be difficult, but it is still far easier than using its name, or attempting to guess at the structure and functionality of a molecule using just its chemical formula of codice_2. Examples of Organic Chemistry Notation. Typically organic chemistry molecules are drawn as modified Lewis structures. If you remember, a Lewis structure uses lines to connect chemical symbols together, illustrating a covalent bond, and also uses dots to represent non-bond electrons. This is shown in the diagram below of carbon dioxide. The drawing illustrates the four electrons of carbon participating in two double bonds with two oxygen atoms, and also the non-bonding electron pairs for each atom of oxygen. In organic chemistry, there are a lot of carbons in every molecule, generally, so organic chemists by convention do not draw every single carbon in every molecule. The same is true of hydrogens attached to the carbons; it is twice the annoyance to draw thirty hydrogens in a fatty acid than it is to draw the fifteen carbons. Therefore, in organic chemistry, carbon atoms are assumed to be wherever a line or line segment begins or ends. Furthermore, enough hydrogen atoms are assumed to be attached to any carbon not marked with a + or - sign (indicating an ionic charge) to bring that carbon's total number of bonds to four. At first this notation may be confusing, but the shorthand method rapidly proves its worth. 

Computer programming is the craft of writing useful, maintainable, and extensible source code which can be interpreted or compiled by a computing system to perform a meaningful task. Programming a computer can be performed in one of numerous languages, ranging from a higher-level language to writing directly in low-level machine code (that is, code that more directly controls the specifics of the computer's hardware) all the way down to writing microcode (which does directly control the electronics in the computer). Using programming languages and markup languages (such as XHTML and XForms) require some of the same skills, but using markup languages is generally not considered "programming." Nevertheless, many markup languages allow inclusion of scripts, e.g. many HTML documents contain JavaScript. There are exceptions where markup languages do represent programming such as SuperX++ (http://xplusplus.sourceforge.net/) and o:XML (http://www.o-xml.org/) Computer programming is one part of a much larger discipline known as software engineering, which includes several different aspects of making software including design, construction and quality control. The subject of this book is software construction, that is, programming. Computer programming is also a useful skill (though not always necessary) for people who are interested in . Whereas software engineering is interested specifically in making software, computer science tends to be oriented towards more theoretical or mathematical problems. Getting started. Many people think they must choose a specific programming language in order to become a programmer, believing that they can only do that language. They ask themselves, "Should I be a C programmer or a Java programmer?" That's completely the wrong question. The right question is "How can I become a good programmer?" Unfortunately the employment market has contributed greatly to misconceptions about computer programming by companies advertising for employees with a specific (therefore limited) computer language skill-set and responses being handled by human resources(HR), without someone with a programming background. There are a few points one can make about what a good programmer knows about specific computer languages. First - many languages are based on the same fundamental building blocks. Learning a language should be seen more as a way of acquiring those concepts than language or machine specific techniques. Second - good programmers are generally competent in more than one language because it is naturally interesting and useful to find different ways of solving problems. It is not necessary to master many different languages or even more than one—a programmer could excel in one language and have only a vague working idea how to program others. It is useful to know many different methods for solving computer problems, also known as algorithms. An algorithm is a list of well-defined instructions for completing a task, and knowing several languages means having the ability to list the computer instructions in many different ways. Since computer programming languages have so much in common, it is generally easy to learn a new programming language once you have mastered another. So how do you get started? One reasonable technique would be to just pick a language and run with it. Unfortunately, we cannot suggest what the right computer language might be for all people for all purposes. Ask ten programmers what language you should learn and you will get ten different responses. Given the collaborative nature of this wikibook, you'll probably get as many responses as there are programming language books on the site. Families of languages. There is a common misconception by people unfamiliar with computer programming that all programming languages are essentially the same. In one sense this is true because all digital electronic computers translate programming languages into strings of ones and zeros called binary, or Machine code. While mainstream, personal computer languages tend to be derived from a specific tradition and are very similar (hence the popularity of this misconception), some languages fall into different paradigms which provide for a radically different programming experience. Programming in Java is quite different from programming in , which is quite different from programming in Haskell or Prolog or Forth, etc. In the American Scientist article The Semicolon Wars, Brian Hayes classifies languages into four categories: imperative, object-oriented, functional, and declarative. Imperative and object-oriented languages tend to be used in the mainstream, whereas functional and declarative languages tend to be used in academic settings. Functional and declarative programming enthusiasts might argue that the paradigms are 20 years ahead of the mainstream and superior in many respects; however, mainstream language advocates would probably counter that such paradigms are hard to learn, or not very practical for their own unpopularity, among other things. We do not make any claims about who is right on this matter, but at the very least, we will suggest that building familiarity with the four major paradigms is an extremely valuable exercise. When it comes to computers, all things are made, and function primarily by, programming. Although programming is an essential part of the functionality of any computer or application, not all programming languages are the same. In fact, they are very different from one another with different uses, functionality, and different levels of complexity. A programming language, in the most basic way, is a set of rules or guidelines that is used to write the computer programs. Even though you are writing the program, you may need a certain type of software or program for the language that you use. There are many different types of programming languages that can be used and each has a different set of rules. Programming has two basic categories. There are low-level and high-level languages, the difference between the two is that low-level languages often use 0s and 1s, and this works because it gives the computer the ability to quickly understand what needs to be done or executed. High level languages are easier to write because they are much closer to the English language and are much more flexible to write with, although there are also different levels of this readability as well and different categories of these languages that can be written. A few examples would be Visual Basic, C++, and Java. Common concepts. Programming languages tend to have many general concepts in common. One can examine the recurring concepts and how they are expressed in various languages in the following table. To see a comparison of syntax in various programming languages, see these "Hello World" examples. For a list including various computer languages arranged together by syntax terms and patterns, see . Programming skills. Computer programming is really just about solving problems. It turns out that a large number of the problems you encounter in the real world are really just special cases of a more general problem. Luckily for you, many of these problems have been studied by computer scientists for a very long time, sometimes leading to provably unbeatable solutions, or sometimes solutions which are "good enough" for every day needs. In short, learning a language gives you skills, but learning data structures and algorithms shows you how to use these skills wisely. History of programming. Specific languages. The following languages deserve special mention, being significant languages in the development of structured programming languages and object-oriented programming. They are worth understanding for the concepts they introduced. Additional Information. Editors. An editor is simply a text program like Notepad. It is said that a programmer's best friend is the editor. A good editor is lightweight, has only essential tools and should support syntax highlighting for your language. Examples of good editors for which we have teaching books are : For more text editors, see Wikipedia's . Popular libraries. Unix native Windows "native" Cross platform 

Like most programming languages, C uses and processes variables. In C, variables are human-readable names for the computer's memory addresses used by a running program. Variables make it easier to store, read and change the data within the computer's memory by allowing you to associate easy-to-remember labels for the memory addresses that store your program's data. The memory addresses associated with variables aren't determined until after the program is compiled and running on the computer. At first, it's easiest to imagine variables as placeholders for values, much like in mathematics. You can think of a variable as being equivalent to its assigned value. So, if you have a variable "i" that is initialized (set equal) to 4, then it follows that "i + 1" will equal "5". However, a skilled C programmer is more mindful of the invisible layer of abstraction going on just under the hood: that a variable is a stand-in for the memory address where the data can be found, not the data itself. You will gain more clarity on this point when you learn about pointers. Since C is a relatively low-level programming language, before a C program can utilize memory to store a variable it must claim the memory needed to store the values for a variable. This is done by declaring variables. Declaring variables is the way in which a C program shows the number of variables it needs, what they are going to be named, and how much memory they will need. Within the C programming language, when managing and working with variables, it is important to know the "type" of variables and the "size" of these types. A type’s size is the amount of computer memory required to store one value of this type. Since C is a fairly low-level programming language, the size of types can be specific to the hardware and compiler used – that is, how the language is made to work on one type of machine can be different from how it is made to work on another. All variables in C are typed. That is, every variable declared must be assigned as a certain type of variable. Declaring, Initializing, and Assigning Variables. Here is an example of declaring an integer, which we've called some_number. (Note the semicolon at the end of the line; that is how your compiler separates one program "statement" from another.) int some_number; This statement tells the compiler to create a variable called codice_1 and associate it with a memory location on the computer. We also tell the compiler the type of data that will be stored at that address, in this case an integer. Note that in C we must specify the type of data that a variable will store. This lets the compiler know how much total memory to set aside for the data (on most modern machines an codice_2 is 4 bytes in length). We'll look at other data types in the next section. Multiple variables can be declared with one statement, like this: int anumber, anothernumber, yetanothernumber; In early versions of C, variables had to be declared at the beginning of a block. In C99 it is allowed to mix declarations and statements arbitrarily – but doing so is not usual, because it is rarely necessary, some compilers still don’t support C99 (portability), and it may, because it is uncommon yet, irritate fellow programmers (maintainability). After declaring variables, you can assign a value to a variable later on using a statement like this: some_number = 3; The assignment of a value to a variable is called "initialization". The above statement directs the compiler to insert an integer representation of the number "3" into the memory address associated with codice_1. We can save a bit of typing by declaring "and" assigning data to a memory address at the same time: int some_new_number = 4; You can also assign variables to the value of other variable, like so: some_number = some_new_number; Or assign multiple variables the same value with one statement: anumber = anothernumber = yetanothernumber = 8; This is because the assignment x = y returns the value of the assignment, y. For example, some_number = 4 returns 4. That said, x = y = z is really a shorthand for x = (y = z). Naming Variables. Variable names in C are made up of letters (upper and lower case) and digits. The underscore character ("_") is also permitted. Names must not begin with a digit. Unlike some languages (such as Perl and some BASIC dialects), C does not use any special prefix characters on variable names. Some examples of valid (but not very descriptive) C variable names: foo Bar BAZ foo_bar _foo42 _ QuUx Some examples of invalid C variable names: 2foo (must not begin with a digit) my foo (spaces not allowed in names) $foo ($ not allowed -- only letters, and _) while (language keywords cannot be used as names) As the last example suggests, certain words are reserved as keywords in the language, and these cannot be used as variable names. It is not allowed to use the same name for multiple variables in the same scope. When working with other developers, you should therefore take steps to avoid using the same name for global variables or function names. Some large projects adhere to naming guidelines to avoid duplicate names and for consistency. In addition there are certain sets of names that, while not language keywords, are reserved for one reason or another. For example, a C compiler might use certain names "behind the scenes", and this might cause problems for a program that attempts to use them. Also, some names are reserved for possible future use in the C standard library. The rules for determining exactly what names are reserved (and in what contexts they are reserved) are too complicated to describe here, and as a beginner you don't need to worry about them much anyway. For now, just avoid using names that begin with an underscore character. The naming rules for C variables also apply to naming other language constructs such as function names, struct tags, and macros, all of which will be covered later. Literals. Anytime within a program in which you specify a value explicitly instead of referring to a variable or some other form of data, that value is referred to as a literal. In the initialization example above, 3 is a literal. Literals can either take a form defined by their type (more on that soon), or one can use hexadecimal (hex) notation to directly insert data into a variable regardless of its type. Hex numbers are always preceded with "0x". For now, though, you probably shouldn't be too concerned with hex. The Four Basic Data Types. In Standard C there are four basic data types. They are codice_2, codice_5, codice_6, and codice_7. The codice_2 type. The int type stores integers in the form of "whole numbers". An integer is typically the size of one machine word, which on most modern home PCs is 32 bits (4 octets). Examples of literals are whole numbers (integers) such as 1, 2, 3, 10, 100... When int is 32 bits (4 octets), it can store any whole number (integer) between -2147483648 and 2147483647. A 32 bit word (number) has the possibility of representing any one number out of 4294967296 possibilities (2 to the power of 32). If you want to declare a new int variable, use the int keyword. For example: int numberOfStudents, i, j=5; In this declaration we declare 3 variables, numberOfStudents, i and j, j here is assigned the literal 5. The codice_5 type. The codice_5 type is capable of holding any member of the execution character set. It stores the same kind of data as an codice_2 (i.e. integers), but typically has a size of one byte. The size of a byte is specified by the macro codice_12 which specifies the number of bits in a char (byte). In standard C it never can be less than 8 bits. A variable of type codice_5 is most often used to store character data, hence its name. Most implementations use the ASCII character set as the execution character set, but it's best not to know or care about that unless the actual values are important. Examples of character literals are 'a', 'b', '1', etc., as well as some special characters such as 'codice_14' (the null character) and 'codice_15' (newline, recall "Hello, World"). Note that the codice_5 value must be enclosed within single quotations. When we initialize a character variable, we can do it two ways. One is preferred, the other way is bad programming practice. The first way is to write char letter1 = 'a'; This is "good" programming practice in that it allows a person reading your code to understand that letter1 is being initialized with the letter 'a' to start off with. The second way, which should "not" be used when you are coding letter characters, is to write char letter2 = 97; /* in ASCII, 97 = 'a' */ This is considered by some to be extremely bad practice, if we are using it to store a character, not a small number, in that if someone reads your code, most readers are forced to look up what character corresponds with the number 97 in the encoding scheme. In the end, codice_17 and codice_18 store both the same thing – the letter 'a', but the first method is clearer, easier to debug, and much more straightforward. One important thing to mention is that characters for numerals are represented differently from their corresponding number, i.e. '1' is not equal to 1. In short, any single entry that is enclosed within 'single quotes'. There is one more kind of literal that needs to be explained in connection with chars: the string literal. A string is a series of characters, usually intended to be displayed. They are surrounded by double quotations (" ", not ' '). An example of a string literal is the "Hello, World!\n" in the "Hello, World" example. The string literal is assigned to a character array, arrays are described later. Example: const char MY_CONSTANT_PEDANTIC_ITCH[] = "learn the usage context.\n"; printf("Square brackets after a variable name means it is a pointer to a string of memory blocks the size of the type of the array element.\n"); The codice_6 type. codice_6 is short for floating point. It stores inexact representations of real numbers, both integer and non-integer values. It can be used with numbers that are much greater than the greatest possible codice_2. codice_6 literals must be suffixed with F or f. Examples are: 3.1415926f, 4.0f, 6.022e+23f. It is important to note that floating-point numbers are inexact. Some numbers like 0.1f cannot be represented exactly as codice_6s but will have a small error. Very large and very small numbers will have less precision and arithmetic operations are sometimes not associative or distributive because of a lack of precision. Nonetheless, floating-point numbers are most commonly used for approximating real numbers and operations on them are efficient on modern microprocessors. Floating-point arithmetic is explained in more detail on Wikipedia. codice_6 variables can be declared using the float keyword. A codice_6 is only one machine word in size. Therefore, it is used when less precision than a double provides is required. The codice_7 type. The double and float types are very similar. The float type allows you to store single-precision floating point numbers, while the double keyword allows you to store double-precision floating point numbers – real numbers, in other words. Its size is typically two machine words, or 8 bytes on most machines. Examples of double literals are 3.1415926535897932, 4.0, 6.022e+23 (scientific notation). If you use 4 instead of 4.0, the 4 will be interpreted as an int. The distinction between floats and doubles was made because of the differing sizes of the two types. When C was first used, space was at a minimum and so the judicious use of a float instead of a double saved some memory. Nowadays, with memory more freely available, you rarely need to conserve memory like this – it may be better to use doubles consistently. Indeed, some C implementations use doubles instead of floats when you declare a float variable. If you want to use a double variable, use the double keyword. sizeof. If you have any doubts as to the amount of memory actually used by any variable (and this goes for types we'll discuss later, also), you can use the sizeof operator to find out for sure. (For completeness, it is important to mention that sizeof is a unary operator, not a function.) Its syntax is: sizeof object sizeof(type) The two expressions above return the size of the object and type specified, in bytes. The return type is size_t (defined in the header &lt;stddef.h&gt;) which is an unsigned value. Here's an example usage: size_t size; int i; size = sizeof(i); size will be set to 4, assuming CHAR_BIT is defined as 8, and an integer is 32 bits wide. The value of sizeof's result is the number of bytes. Note that when sizeof is applied to a char, the result is 1; that is: sizeof(char) always returns 1. Data type modifiers. One can alter the data storage of any data type by preceding it with certain modifiers. long and short are modifiers that make it possible for a data type to use either more or less memory. The int keyword need not follow the short and long keywords. This is most commonly the case. A short can be used where the values fall within a lesser range than that of an int, typically -32768 to 32767. A long can be used to contain an extended range of values. It is not guaranteed that a short uses less memory than an int, nor is it guaranteed that a long takes up more memory than an int. It is only guaranteed that sizeof(short) &lt;= sizeof(int) &lt;= sizeof(long). Typically a short is 2 bytes, an int is 4 bytes, and a long either 4 or 8 bytes. Modern C compilers also provide long long which is typically an 8 byte integer. In all of the types described above, one bit is used to indicate the sign (positive or negative) of a value. If you decide that a variable will never hold a negative value, you may use the unsigned modifier to use that one bit for storing other data, effectively doubling the range of values while mandating that those values be positive. The unsigned specifier also may be used without a trailing int, in which case the size defaults to that of an int. There is also a signed modifier which is the opposite, but it is not necessary, except for certain uses of char, and seldom used since all types (except char) are signed by default. The long modifier can also be used with double to create a long double type. This floating-point type may (but is not required to) have greater precision than the double type. To use a modifier, just declare a variable with the data type and relevant modifiers: unsigned short int usi; /* fully qualified -- unsigned short int */ short si; /* short int */ unsigned long uli; /* unsigned long int */ const qualifier. When the const qualifier is used, the declared variable must be initialized at declaration. It is then not allowed to be changed. While the idea of a variable that never changes may not seem useful, there are good reasons to use const. For one thing, many compilers can perform some small optimizations on data when it knows that data will never change. For example, if you need the value of π in your calculations, you can declare a const variable of pi, so a program or another function written by someone else cannot change the value of pi. Note that a Standard conforming compiler must issue a warning if an attempt is made to change a const variable - but after doing so the compiler is free to ignore the const qualifier. Magic numbers. When you write C programs, you may be tempted to write code that will depend on certain numbers. For example, you may be writing a program for a grocery store. This complex program has thousands upon thousands of lines of code. The programmer decides to represent the cost of a can of corn, currently 99 cents, as a literal throughout the code. Now, assume the cost of a can of corn changes to 89 cents. The programmer must now go in and manually change each entry of 99 cents to 89. While this is not that big a problem, considering the "global find-replace" function of many text editors, consider another problem: the cost of a can of green beans is also initially 99 cents. To reliably change the price, you have to look at every occurrence of the number 99. C possesses certain functionality to avoid this. This functionality is approximately equivalent, though one method can be useful in one circumstance, over another. Using the const keyword. The const keyword helps eradicate magic numbers. By declaring a variable const corn at the beginning of a block, a programmer can simply change that const and not have to worry about setting the value elsewhere. There is also another method for avoiding magic numbers. It is much more flexible than const, and also much more problematic in many ways. It also involves the preprocessor, as opposed to the compiler. Behold... #define. When you write programs, you can create what is known as a "macro", so when the computer is reading your code, it will replace all instances of a word with the specified expression. Here's an example. If you write when you want to, for example, print the price of corn, you use the word codice_27 instead of the number 0.99 – the preprocessor will replace all instances of codice_27 with 0.99, which the compiler will interpret as the literal codice_7 0.99. The preprocessor performs substitution, that is, codice_27 is replaced by 0.99 so this means there is no need for a semicolon. It is important to note that codice_31 has basically the same functionality as the "find-and-replace" function in a lot of text editors/word processors. For some purposes, codice_31 can be harmfully used, and it is usually preferable to use codice_33 if codice_31 is unnecessary. It is possible, for instance, to codice_31, say, a macro codice_36 as the number 3, but if you try to print the macro, thinking that codice_36 represents a string that you can show on the screen, the program will have an error. codice_31 also has no regard for type. It disregards the structure of your program, replacing the text "everywhere" (in effect, disregarding scope), which could be advantageous in some circumstances, but can be the source of problematic bugs. You will see further instances of the codice_31 directive later in the text. It is good convention to write codice_31d words in all capitals, so a programmer will know that this is not a variable that you have declared but a codice_31d macro. It is not necessary to end a preprocessor directive such as codice_31 with a semicolon; in fact, some compilers may warn you about unnecessary tokens in your code if you do. Scope. In the Basic Concepts section, the concept of scope was introduced. It is important to revisit the distinction between local types and global types, and how to declare variables of each. To declare a local variable, you place the declaration at the beginning (i.e. before any non-declarative statements) of the block to which the variable is deemed to be local. To declare a global variable, declare the variable outside of any block. If a variable is global, it can be read, and written, from anywhere in your program. Global variables are not considered good programming practice, and should be avoided whenever possible. They inhibit code readability, create naming conflicts, waste memory, and can create difficult-to-trace bugs. Excessive usage of globals is usually a sign of laziness or poor design. However, if there is a situation where local variables may create more obtuse and unreadable code, there's no shame in using globals. Other Modifiers. Included here, for completeness, are more of the modifiers that standard C provides. For the beginning programmer, "static" and "extern" may be useful. "volatile" is more of interest to advanced programmers. "register" and "auto" are largely deprecated and are generally not of interest to either beginning or advanced programmers. static. static is sometimes a useful keyword. It is a common misbelief that the only purpose is to make a variable stay in memory. When you declare a function or global variable as "static", you cannot access the function or variable through the extern (see below) keyword from other files in your project. This is called "static linkage". When you declare a local variable as "static", it is created just like any other variable. However, when the variable goes out of scope (i.e. the block it was local to is finished) the variable stays in memory, retaining its value. The variable stays in memory until the program ends. While this behaviour resembles that of global variables, static variables still obey scope rules and therefore cannot be accessed outside of their scope. This is called "static storage duration". Variables declared static are initialized to zero (or for pointers, NULL) by default. They can be initialized explicitly on declaration to any "constant" value. The initialization is made just once, at compile time. You can use static in (at least) two different ways. Consider this code, and imagine it is in a file called jfile.c: static int j = 0; void up(void)  /* k is set to 0 when the program starts. The line is then "ignored"  * for the rest of the program (i.e. k is not set to 0 every time up()  * is called)  static int k = 0;  j++;  k++;  printf("up() called. k= %2d, j= %2d\n", k , j);  void down(void)  static int k = 0;  j--;  k--;  printf("down() called. k= %2d, j= %2d\n", k , j);  int main(void)  int i;  /* call the up function 3 times, then the down function 2 times */  for (i = 0; i &lt; 3; i++)  up();  for (i = 0; i &lt; 2; i++)  down();  return 0; The codice_43 variable is accessible by both up and down and retains its value. The codice_44 variables also retain their value, but they are two different variables, one in each of their scopes. Static variables are a good way to implement encapsulation, a term from the object-oriented way of thinking that effectively means not allowing changes to be made to a variable except through function calls. Running the program above will produce the following output: up() called. k= 1, j= 1 up() called. k= 2, j= 2 up() called. k= 3, j= 3 down() called. k= -1, j= 2 down() called. k= -2, j= 1 Features of codice_45 variables :  1. Keyword used - static  2. Storage - Memory  3. Default value - Zero  4. Scope - Local to the block in which it is declared  5. Lifetime - Value persists between different function calls  6. Keyword optionality - Mandatory to use the keyword extern. extern is used when a file needs to access a variable in another file that it may not have #included directly. Therefore, "extern" does not allocate memory for the new variable, it just provides the compiler with sufficient information to access a variable declared in another file. Features of codice_46 variable :  1. Keyword used - extern  2. Storage - Memory  3. Default value - Zero  4. Scope - Global (all over the program)  5. Lifetime - Value persists till the program's execution comes to an end  6. Keyword optionality - Optional if declared outside all the functions volatile. volatile is a special type of modifier which informs the compiler that the value of the variable may be changed by external entities other than the program itself. This is necessary for certain programs compiled with optimizations – if a variable were not defined volatile then the compiler may assume that certain operations involving the variable are safe to optimize away when in fact they aren't. "volatile" is particularly relevant when working with embedded systems (where a program may not have complete control of a variable) and multi-threaded applications. auto. auto is a modifier which specifies an "automatic" variable that is automatically created when in scope and destroyed when out of scope. If you think this sounds like pretty much what you've been doing all along when you declare a variable, you're right: all declared items within a block are implicitly "automatic". For this reason, the "auto" keyword is more like the answer to a trivia question than a useful modifier, and there are lots of very competent programmers that are unaware of its existence. Features of codice_47 variables :  1. Keyword used - auto  2. Storage - Memory  3. Default value - Garbage value (random value)  4. Scope - Local to the block in which it is defined  5. Lifetime - Value persists while the control remains within the block  6. Keyword optionality - Optional register. register is a hint to the compiler to attempt to optimize the storage of the given variable by storing it in a register of the computer's CPU when the program is run. Most optimizing compilers do this anyway, so use of this keyword is often unnecessary. In fact, ANSI C states that a compiler can ignore this keyword if it so desires – and many do. Microsoft Visual C++ is an example of an implementation that completely ignores the "register" keyword. Features of codice_48 variables :  1. Keyword used - register  2. Storage - CPU registers (values can be retrieved faster than from memory)  3. Default value - Garbage value  4. Scope - Local to the block in which it is defined  5. Lifetime - Value persists while the control remains within the block  6. Keyword optionality - Mandatory to use the keyword 

An assignment statements in wikibook pseudocode is written as codice_1.  let X := 10 But this is not needed in wikibook pseudocode. 

Welcome. Welcome to English in Use, a book about the actual use of the English language. It features sections about types of grammar, punctuation and formality. This wikibook is intended for use by native speakers of English or advanced learners of English as a second/foreign language. If you wish to learn English, then you should use one of the English books for students learning English as a second/foreign language. For new learners of English, try English for beginners book. If you have some experience of English but need to refresh your knowledge, try the English for B2 Students Wikibook. If you're preparing for the Cambridge FCE exam (Cambridge English: First), then head to the module. Those interested in English for business contexts should try the Business English Wikibook. 

Computer Programming Types determine the kinds of values and how they can be used in the given programming environment. In most cases, a programming language defines a set of basic data types, e.g. for numbers, characters or strings. In higher languages, it is often possible to define new data types from the existing ones, for example, to represent a postal address (consisting of strings for street and city and integers for the postal code). When communicating data between different programs and computer systems it is important to either use types that both can recognize, or have a means of translating between them. Type Definition. A complete definition of a data type consists of up to three things: Primitive Types. There are a few basic data types seen in some programming languages: 

Arrays. An array is a collection, mainly of similar data types, stored into a common variable. The collection forms a data structure where objects are stored linearly, one after another in memory. Sometimes arrays are even replicated into the memory hardware. The structure can also be defined as a particular method of storing elements of indexed data. Elements of data are logically stored sequentially in blocks within the array. Each element is referenced by an index, or subscripts. The index is usually a number used to address an element in the array. For example, if you were storing information about each day in August, you would create an array with an index capable of addressing 31 values—one for each day of the month. Indexing rules are language dependent, however most languages use either 0 or 1 as the first element of an array. The concept of an array can be daunting to the uninitiated, but it is really quite simple. Think of a notebook with pages numbered 1 through 12. Each page may or may not contain information on it. The notebook is an "array" of pages. Each page is an "element" of the array 'notebook'. Programmatically, you would retrieve information from a page by referring to its number or "subscript", i.e., notebook(4) would refer to the contents of page 4 of the array notebook. &lt;br&gt;"The notebook (array) contains 12 pages (elements)" Arrays can also be multidimensional - instead of accessing an element of a one-dimensional list, elements are accessed by two or more indices, as from a matrix or tensor. Multidimensional arrays are as simple as our notebook example above. To envision a multidimensional array, think of a calendar. Each page of the calendar, 1 through 12, is an element, representing a month, which contains approximately 30 elements, which represent days. Each day may or may not have information in it. Programmatically then, calendar(4,15) would refer to the 4th month, 15th day. Thus we have a two-dimensional array. To envision a three-dimensional array, break each day up into 24 hours. Now calendar(4,15,9) would refer to 4th month, 15th day, 9th hour. &lt;br&gt;"A simple 6 element by 4 element array" Arrays guarantee constant time read and write access, formula_1, however many lookup operations (find_min, find_max, find_index) of an instance of an element are linear time, formula_2. Arrays are very efficient in most languages, as operations compute the address of an element via a simple formula based on the base address element of the array. Array implementations differ greatly between languages: some languages allow arrays to be re-sized automatically, or to even contain elements of differing types (such as Perl). Other languages are very strict and require the type and length information of an array to be known at run time (such as C). Arrays typically map directly to contiguous storage locations within your computer's memory and are therefore the "natural" storage structure for most higher level languages. Simple linear arrays are the basis for most of the other data structures. Many languages do not allow you to allocate any structure except an array, everything else must be implemented on top of the array. The exception is the linked list, that is typically implemented as individually allocated objects, but it is possible to implement a linked list within an array. Type. The array index needs to be of some type. Usually, the standard integer type of that language is used, but there are also languages such as Ada and Pascal which allow any discrete type as an array index. Scripting languages often allow any type as an index (associative array). Bounds. The array index consists of a range of values with a lower bound and an upper bound. In some programming languages only the upper bound can be chosen while the lower bound is fixed to be either 0 or 1 . In other programming languages both the upper and lower bound can be freely chosen . Bounds check. The third aspect of an array index is the check for valid ranges and what happens when an invalid index is accessed. This is a very important point since the majority of computer worms and computer viruses attack by using invalid array bounds. There are three options open: Declaring Array Types. The declaration of array type depends on how many features the array in a particular language has. The easiest declaration is when the language has a fixed lower bound and fixed index type. If you need an array to store the monthly income you could declare in C typedef double Income[12]; This gives you an array with in the range of 0 to 11. For a full description of arrays in C see C Programming/Arrays. If you use a language where you can choose both the lower bound as well as the index type, the declaration is—of course—more complex. Here are two examples in Ada: type Month is range 1 .. 12; type Income is array(Month) of Float; or shorter: type Income is array(1 .. 12) of Float; For a full description of arrays in Ada see Ada Programming/Types/array. Array Access. We generally write arrays with a name, followed by the index in some brackets, square '[]' or round '()'. For example, August[3] is the method used in the C programming language to refer to a particular day in the month. Because the C language starts the index at zero, August[3] is the 4th element in the array. august[0] actually refers to the first element of this array. Starting an index at zero is natural for computers, whose internal representations of numbers begin with zero, but for humans, this unnatural numbering system can lead to problems when accessing data in an array. When fetching an element in a language with zero-based indexes, keep in mind the "true" length of an array, lest you find yourself fetching the wrong data. This is the disadvantage of programming in languages with fixed lower bounds, the programmer must always remember that "[0]" means "1st" and, when appropriate, add or subtract one from the index. Languages with variable lower bounds will take that burden off the programmer's shoulder. We use indexes to store "related" data. If our C language array is called august, and we wish to store that we're going to the supermarket on the 1st, we can say, for example  august[0] = "Going to the shops today" In this way, we can go through the indexes from 0 to 30 and get the related tasks for each day in august. 



Grammar - Preterite (el pretérito indefinido). The following table shows the preterite of regular verbs. Regular -er and -ir verbs follow the exact same pattern. Note that the nosotros form is the same as in the present tense for -ar and -ir verbs, so you have to know the context to figure out the time. Also, note that the last letter of comí and viví has an accent mark. Here is a list of common verbs that have an irregular preterite: Exercise: Preterite Grammar - Imperfect (el pretérito imperfecto). The following table shows the imperfect of regular verbs. Note that regular -er and -ir verbs follow the exact same pattern: There are only three verbs that are irregular in the imperfect: Grammar - Preterite vs. Imperfect. Spanish has two tenses that correspond to the English "simple past". Roughly speaking, the Preterite is used to tell What happened - it refers to a specific event. The Imperfect is used to tell How things were - it refers to the general situation. Exercise: Preterite vs. Imperfect 

Arrays in C act to store related data under a single variable name with an index, also known as a "subscript". It is easiest to think of an array as simply a list or ordered grouping for variables of the same type. As such, arrays often help a programmer organize collections of data efficiently and intuitively. Later we will consider the concept of a "pointer", fundamental to C, which extends the nature of the array (array can be termed as a constant pointer). For now, we will consider just their declaration and their use. Arrays. C arrays are declared in the following form: type name[number of elements]; For example, if we want an array of six integers (or whole numbers), we write in C: int numbers[6]; For a six character array called "letters", char letters[6]; and so on. You can also initialize as you declare. Just put the initial elements in curly brackets separated by commas as the initial value: For example, if we want to initialize an array with six integers, with 0, 0, 1, 0, 0, 0 as the initial values: int point[6]={0,0,1,0,0,0}; Though when the array is initialized as in this case, the array dimension may be omitted, and the array will be automatically sized to hold the initial data: int point[]={0,0,1,0,0,0}; This is very useful in that the size of the array can be controlled by simply adding or removing initializer elements from the definition without the need to adjust the dimension. If the dimension is specified, but not all elements in the array are initialized, the remaining elements will contain a value of 0. This is very useful, especially when we have very large arrays. int numbers[2000]={245}; The above example sets the first value of the array to 245, and the rest to 0. If we want to access a variable stored in an array, for example with the above declaration, the following code will store a 1 in the variable x int x; x = point[2]; Arrays in C are indexed starting at 0, as opposed to starting at 1. The first element of the array above is point[0]. The index to the last value in the array is the array size minus one. In the example above the subscripts run from 0 through 5. C does not guarantee bounds checking on array accesses. The compiler may not complain about the following (though the best compilers do): char y; int z = 9; char point[6] = { 1, 2, 3, 4, 5, 6 }; //examples of accessing outside the array. A compile error is not always raised y = point[15]; y = point[-4]; y = point[z]; During program execution, an out of bounds array access does not always cause a run time error. Your program may happily continue after retrieving a value from point[-1]. To alleviate indexing problems, the sizeof() expression is commonly used when coding loops that process arrays. Many people use a macro that in turn uses sizeof() to find the number of elements in an array, a macro variously named "lengthof()", "MY_ARRAY_SIZE()" or "NUM_ELEM()", "SIZEOF_STATIC_ARRAY()", etc. int ix; short anArray[]= { 3, 6, 9, 12, 15 }; for (ix=0; ix&lt; (sizeof(anArray)/sizeof(short)); ++ix) {  DoSomethingWith("%d", anArray[ix] ); Notice in the above example, the size of the array was not explicitly specified. The compiler knows to size it at 5 because of the five values in the initializer list. Adding an additional value to the list will cause it to be sized to six, and because of the sizeof expression in the for loop, the code automatically adjusts to this change. Good programming practice is to declare a variable size , and store the number of elements in the array in it. size = sizeof(anArray)/sizeof(short) C also supports multi dimensional arrays (or, rather, arrays of arrays). The simplest type is a two dimensional array. This creates a rectangular array - each row has the same number of columns. To get a char array with 3 rows and 5 columns we write in C  char two_d[3][5]; To access/modify a value in this array we need two subscripts: char ch; ch = two_d[2][4]; or two_d[0][0] = 'x'; Similarly, a multi-dimensional array can be initialized like this: int two_d[2][3] = ; The number of columns must be explicitly stated; however, the compiler will find the appropriate amount of rows based on the initializer list. There are also weird notations possible: int a[100]; int i = 0; if (a[i]==i[a])  printf("Hello world!\n"); a[i] and i[a] refer to the same location. (This is explained later in the next Chapter.) Strings. C has no string handling facilities built in; consequently, strings are defined as arrays of characters. C allows a character array to be represented by a character string rather than a list of characters, with the null terminating character automatically added to the end. For example, to store the string "Merkkijono", we would write char string[11] = "Merkkijono"; or char string[11] = {'M', 'e', 'r', 'k', 'k', 'i', 'j', 'o', 'n', 'o', '\0'}; In the first example, the string will have a null character automatically appended to the end by the compiler; by convention, library functions expect strings to be terminated by a null character. The latter declaration indicates individual elements, and as such the null terminator needs to be added manually. Strings do not always have to be linked to an explicit variable. As you have seen already, a string of characters can be created directly as an unnamed string that is used directly (as with the printf functions.) To create an extra long string, you will have to split the string into multiple sections, by closing the first section with a quote, and recommencing the string on the next line (also starting and ending in a quote): char string[58] = "This is a very, very long "  "string that requires two lines."; While strings may also span multiple lines by putting the backslash character at the end of the line, this method is deprecated. There is a useful library of string handling routines which you can use by including another header file. This standard string library will allow various tasks to be performed on strings, and is discussed in the Strings chapter. 

Welcome to the Wikibook of Calculus This wikibook aims to be a high quality calculus textbook through which users can master the discipline. Standard topics such as "limits", "differentiation" and "integration" are covered, as well as several others. Please contribute wherever you feel the need. You can simply help by rating individual sections of the book that you feel were inappropriately rated! 

Kanji () characters are based on Chinese characters transmitted to Japan during the spread of Buddhism in the 5th century. A large percentage (approx. 70%) of Japanese vocabulary comes from Chinese or Chinese-derived words. While the meaning of individual characters is fairly consistent between the languages, compound words often have different meanings. Kanji are inflected by hiragana that follow and particles give the case. Most words are written using kanji, though some have none and loan-words from other languages are generally written in katakana. The large number of homophones makes it highly desirable to use kanji and knowing them can help with memorising new words. Note that "writing" kanji skillfully is significantly harder than "reading" kanji skillfully, since one must "recall" characters, not simply "recognize" them. Further, with Input Methods allowing one to write Japanese on a computer phonetically (by recognizing the kanji, not needing to produce them), the practical need for kanji writing skills is lower than in the past, but it is still fundamental to mastery of Japanese. Study methods. Kanji can form a difficult hurdle for some in their study of Japanese. Their nature as graphic representations of concepts translating to sounds gives rise to the particularly diverse methods employed for the study of kanji. Fundamentally, one’s goal is to learn "Japanese," not "kanji" per se and this has two main implications. Firstly, as many words are written as compounds of multiple kanji it is not sufficient learn the individual two thousand odd characters, but also their combinations. Furthermore, just as learning vocabulary in any language, these must be learnt in the context of the language. Not only does it aid memorisation of terms, but enforces the understanding of their nuance. It is finally worth mentioning that one can learn to speak Japanese without learning to read or write it, just as with any language. If one is, however, ever to learn to read, it is advisable to start right away and learn the characters in parallel with vocabulary and phrases. Throughout, understand that one’s mastery of any skill is imperfect, impermanent, and incomplete (see 侘寂, "wabi-sabi") – while perfection is a worthy goal, it should not be expected nor demanded – mistakes should be expected, and accept that there are further levels of mastery: not 300 people a year pass the Kanji kentei level 1. Basic issues, regardless of study methods: Because there are so many kanji, and they are relatively sparse (of 1,945 kanji, most will not be used and reinforced in any sample of text, unlike kana) simply memorizing the forms and pronunciations (as one does for the 26 letters of the Latin alphabet or the 46 kana, twice) is less practical and effective, and one instead uses more structured mnemonic methods. There are three aspects to a particular kanji: There are a number of ways to learn the kanji. Rather than pick one, try to see how each of these works for you and combine them in your study. Rote. The most straight forward way of learning kanji is by rote. While few will succeed in retaining even a portion of the two thousand basic characters — not to mention their compounds — rote learning is a good way to practice mnemonic devices such as those mentioned in the following sections. Writing reinforces character details, builds muscle memory and improves handwriting. Thus, regardless of learning system, practicing "writing" the kanji is a valuable aspect of learning. Make flash-cards with one or more characters on one side, the meaning and reading on the other and drill yourself. Make another set of cards with the meaning on one side and the characters and readings on the other and drill yourself on writing the kanji. There are several programs and website applications that offer kanji drilling. Notable spaced recognition software include Anki and Mnemosyne. Forgetting a rarely used kanji is easy so it is important to regularly review these. Handwriting. Note that characters have a generally accepted stroke order. While this may seem an extra burden at first, the order is highly regular and will vastly improve your ability to read other people's handwriting, not to mention make yours more intelligible. As with the handwriting of most scripts, Japanese calligraphy has a long history and is greatly revered to this day. As kanji are somewhat more intricate than Latin characters, the quality of handwriting and the order the strokes are written in matter a great deal. In fact, without a commonly accepted system, cursive styles and hurried handwriting would be illegible, indeed. There is, of course, only one way to practice handwriting: By writing. Get yourself a nice notebook, preferable one with good sized squares, and practice, practice, practice. Context. The "Kanji in Context" texts from the Inter-University Center for Japanese Language Studies emphasizes the value of learning not so much kanji characters in isolation, but "kanji-based vocabulary," particularly as part of phrases or idioms. In this approach, when learning a kanji, one learns important words that it is part of. Further, one will learn kanji that make up a given word at the same time – for example, one will learn the "word" 日本 (Nihon, Japan) and, at the same time, the characters 日 (nichi, ni, sun) and 本 (hon, root). Recognising the constituent parts. As you progress in learning kanji, you'll start to see patterns emerge; constituent parts of characters that are common among many characters. Recognising these will allow you to see the characters as made up of shapes rather than just strokes and thus simplify retaining them. The general method is to systematically break up characters into graphical components, some of which may not be used as separate characters. Next, one systematically maps these elements to some mnemonic, and then builds a picture or story combining these. The principles at work are: Useful resources for diagnosing these constituent parts are the book and online version of "Kanji ABC". James Heisig's well-known series "Remembering the Kanji" is the best known study aid that uses this method. Alternatives include Smart Kanji Book which only includes common kanji and the primitives that form them and the Kanji Pict-o-graphix which uses a graphic approach instead of mnemonic stories." A further such resource is Genki’s Kanji Look and Learn. These may not be sufficient in themselves, as they focus purely on the characters, but can be valuable components of one’s learning, helping with remembering character forms and especially minor details. Chinese-derived reading. The vast majority of Chinese characters are composed as "phono-semantic compounds": one component (generally the radical) is semantic (about the meaning), and the other component is used for its phonetic value (sound) called . Note that this is how the character as used in "Chinese" is composed. As kanji usually have several readings, including a Chinese-derived one, this can be used to remember the character and one of its readings. Understanding this, and decomposing characters into Phonetic + Semantic components and relating them to similar characters using either of these components helps with remembering the character’s form, its meaning, and a Chinese-derived pronunciation ("on yomi," 音読み). For example, the character for small is 小 which has the Chinese-derived reading shō. The characters 少, 炒, 抄, 省, 称, 鈔 and 渉 all share that same Chinese reading. Again, keep in mind that these are only the "Chinese" readings and the each of these has other different readings as well. Attention to detail. It is easy to make minor mistakes with both recognizing and writing kanji. A high level of mastery requires attention to detail and continual polish (see 改善, "kaizen"). Even at lower levels, attention to detail yields overlearning and deepens understanding; if you are worrying about the stroke order, you are likely not forgetting the character outright. Minor errors can be made in writing (e.g. incorrect strokes, strokes touching when they should not, or incorrect stroke order) and pronunciation (e.g. incorrect voicing; especially "rendaku"/"euphonic changes"). To achieve a high level requires "detecting" and "correcting" such errors. Realizing that one has forgotten a kanji is easy enough. For other errors, one may not notice them, or one may feel a "lack of confidence" reflecting imperfect mastery. To detect such errors one must review regularly and ensure that all these details are correct. Particularly useful in subtle errors is to study the character in question with various related characters (both graphically, as in , and etymologically), and in the context of various words: this allows one to "contrast" the character, rather than trying to retain it in isolation. Readings. A single Kanji letter can be read (pronounced) in many different ways, depending on its context. These readings are categorized into two main groups - that of Chinese origin (on-yomi, ) and Japanese origin (kun-yomi, ). A third group, the nanoriyomi, is used for the names of people and places. It is often the case that a Kanji letter has more than one reading of Chinese origin. This is because the importing of Chinese letters (with their readings) did not occur just at one time from one region. Onyomi. Onyomi (音読み) is the Chinese-derived reading, which is most commonly used in compound words and for the numbers. It may be useful to note that in most kanji databases, the "on" reading is written in katakana instead of hiragana. 一 (イチ), 二 (ニ), 三 (サン), 四 (シ) are the first four numbers and all are onyomi. Kunyomi. Kunyomi (訓読み) is the Japanese reading, which can be read as a separate word or can be used in compounds. This reading is generally written in hiragana in kanji lists. 月 (つき, tsuki) and 日 (ひ, hi) are the moon and sun and are in kunyomi. Nanoriyomi. Nanoriyomi (名乗り読み) is the name reading, which is used for people's names and for places. Both "康", read as "やす" (e.g. 徳川家康), and "信", read as "のぶ" (e.g. 織田信長), are written in nanoriyomi. Kanji Repetition. The noma: (々), symbol indicates the repetition of a Kanji. The word われわれ indicates "us" or "our group" and is written as "我々" instead of "我我", although they are both the same. The same is true with "人々" (ひとびと), meaning people). JLPT. The Japanese Language Proficiency Test (), or JLPT, is a standardized test to evaluate and certify the language proficiency of non-native Japanese speakers. The JLPT has five levels beginning at level N5 and progressing to level N1 - the most difficult. Each level has a certain set of kanji. JLPT level N5. N5 tests students' recognition of 79 kanji and 482 words. JLPT level N4. N4 tests students' recognition of 166 kanji and 453 words. JLPT level N3. N3 tests students' recognition of 367 kanji and 1555 words. JLPT level N2. N2 tests students' recognition of 367 kanji and 1481 words. JLPT level N1. N1 tests students' recognition of 1231 kanji and 2773 words. 

Otras temas de los verbos. Verbos con la voz pasiva. An action takes place without anyone being assigned responsibility for doing it. Verbos reflexivos. Reflexive verbs indicate that the action of the verb reflects back on the subject. External Links. Spanish Wikibook | ../Verb Tenses/ | ../Verbs List/ 

The Spanish Wikibook was created on August 2, 2003 by ; it was the first language book on Wikibooks. During December 2006, it underwent a complete archive and rewrite, by . A list of major contributors — past and present — can be found below. 

Grammar - Formal Commands (el imperativo). Commands are used when you ask someone to do something or give instructions to people. In this lesson we learn the formal commands, which are the ones you say to persons where you use the usted or ustedes form. The following table shows the endings for the regular verbs. Note that the stem changes that occur in the yo form, (e-&gt;ie, e-&gt;i, o-&gt;ue, ar/er/ir-&gt;go cer-&gt;zco etc., ) apply for formal commands: The following verbs have irregular formal commands: Like in English, the command is usually put in the beginning of the sentence: Examples: Grammar - Informal Tú-Commands. In this lesson we learn the commands you say to someone you would address in the tú form. Spanish distinguishes two different types of tú-commands, the affirmative ("do something") and the negative ("don't do something"). Like the formal commands, we also apply stem changes here: The following verbs have irregular informal tú-commands for the affirmative and negative. &lt;br&gt; This can be memorized with the rhyming mnemonic device "di haz pon ten, sal sé ve ven." Examples: 

The infinitive is the simple, unconjugated form of a verb. It does not indicate who is doing the action or any time reference. "To be" is a verb in the English infinitive form. "Ser" is a verb in the Spanish infinitive form. When you look up a verb in the dictionary it is found in the infinitive form. In Spanish the infinitive forms always end with "ar", "er" or "ir". Examples: "amar" (to love), "temer" (to fear), "partir" (to leave). The infinitive can be used for a lot of things in Spanish grammar. It functions as a gerund in Spanish. In the following examples, I will use "nadar" (to swim): 

Oxidation and reduction. Two important types of reactions in organic chemistry are oxidation and reduction. In oxidation reactions, the oxidized species "loses electron density." In reduction reactions, the reduced species "gains electron density." Of course, these two actions happen in unison as one species is reduced and the other is oxidized. The term redox was coined from the fragments "red" (reduction) and "ox" (oxidation). A standard mnemonic for the terms is “OILRIG”: oxidation is loss, reduction is gain. Oxidation. Oxidation was first observed when oxygen drew electrons off of metals, which were then referred to as "oxidized". (Oxygen is more elecronegative than most other elements.) The term was then applied later to the part of any reaction where electrons are drawn off. Other elements that commonly oxidize in organic reactions include halogens like chlorine and bromine. Reduction. Reduction of a chemical species results in the gain of electrons for that species. This does not necessarily include any change in formal charge; any time an atom increases its electron density even a little bit it is said to be reduced. For example, if an oxygen is removed from a carbon and replaced by a hydrogen (assume the oxygen is also bonded to another atom), the formal charge of the carbon does not change. However, the carbon "sees" a greater share of the electrons from the single bond to hydrogen than it did for the single bond to oxygen. That is because hydrogen is less electronegative than oxygen and gives up its electrons a bit more easily than oxygen does. So a carbon bonded to hydrogen can take up more of its electron density than the same carbon bonded to oxygen. Oxidation numbers. It's possible to assign an “oxidation number” to each atom in a molecule. There are a two different approaches to this. For organic molecules it is generally possible to find all the oxidation numbers using a set of simplified rules. There is no single best set of rules, but as an example, given in order of decreasing priority: From this it is possible to find, for example, that the oxidation state of carbon in methanal: We find that: Denoting the oxidation number of carbon as , we have Carbon has oxidation number zero. Verify for yourself that in methane it has −4, and in ethane −3. An alternative and more general approach is to take the structure of the molecule, and break the bonds such that: As usual, a single bond holds two electrons and a double bond four. (There are further complications which we will not go into here.) The charge of each atom after this process is its oxidation number. Given that oxygen is more electronegative than carbon, which is more electronegative than hydrogen, verify that this gives the same results as before for methanal, methane, and ethane. Both methods give the result that in a neutral species containing only a single element (such as H2 or graphene) all atoms have oxidation state 0. 

Functional groups in reactions make your life easier as an organic chemist because they draw your attention right to where the action is. Any time a reaction is going to occur, you can be almost certain that it is going to take place at a functional group. There are many functional groups of interest to organic chemists. Here are a few: 1. Halides These groups are all made up of a single atom in Group 17 of the Periodic Table, which is known as the halogen group, bonded to a carbon atom. They include fluorine, chlorine, bromine, and iodine. Astatine is also a halogen, but it is rarely discussed because it is not readily found in nature and is radioactive. Their electronegativities vary from fluorine with 4.0 to iodine with 2.5, which is approximately the same value that carbon has. Each of these atoms are able to form a single bond with a carbon atom, replacing hydrogen in alkanes and adding across multiple bonds in alkenes and alkyne. Of the four, fluorine is the most reactive and iodine is the least. Because of their intermediate reactivity, chlorine and bromine are often more useful in many reactions. 2. Carbonyl This group consists of an oxygen atom doubly bonded to a carbon atom. Carbonyl groups are important because the oxygen atom, with an electronegativity of 3.5, shifts electron density away from the rest of the molecule and towards itself. Carbonyls are a key ingredient in aldehydes, ketones, carboxylic acids, esters, and amides. 3. Hydroxyl This group consists of a hydrogen atom singly bonded to an oxygen atom. The electronegativity difference between hydrogen, which has an electronegatity of 2.1, and oxygen, 3.5, pulls electrons away from the hydrogen and makes it somewhat acidic. This acidic character varies depending on the composition of the rest of the molecule. Hydroxyl groups are found in alcohols, phenols, enols, and carboxylic acids. 

Drawing reactions. There are various techniques you will run across for notation of organic reactions. Arrows. Curved arrows are used to show movement of electrons. They are not used to show where atoms, ions, or molecules move, just electrons. A curved arrow with two "hooks" on the end indicated movement of a pair of electrons. A curved arrow with one "hook" indicated movement of one electron. Double-headed arrows are used to represent equivalence between resonance structures. Two-way double half arrows are represent a reaction that can go forward or reverse. If one of the half arrows is longer than the other it means that the reaction pathway favors that direction with the longer arrow. 

&lt;includeonly&gt; Chemical Reactions. &lt;/includeonly&gt; 

The Wikibook Introduction to computers and communications technology. Scope and Audience. This textbook is intended to provide a comprehensive overview of the field, useful to those who will practice in it, and those who need to deal with it peripherally. While it is not a technical work it does deal with technology and describes its structure, background, content, and business in some detail. This means that it while it can be a general survey, it cannot be necessarily simple. Every attempt is made to add no more complex details than are necessary to place a part of the field in context. But context here also means having a general understanding of the internals of an area, and so some complexity is unavoidable. Thus there are three goals established: These goals may at first seem contradictions, but they are not mutually exclusive. Instead they serve as boundaries to identify the material included for each topic. Another limit is imposed with respect to large or "mainframe" computers. The mainframe environment and history are briefly discussed in the introducyion to each part, but most material is based on the now ubiquitous personal or desktop computer environment. Professional Readers. Those who actively work in computers and communications tend to increase the depth of their knowledge and skills within a limited domain. As they accumulate greater experience and skill, their viewpoint tends to become increasingly that of a single role or area within the discipline. While this text can serve them for an initial survey of the field, its greatest value will be as a continued reference and context to provide a higher level view of those areas outside of their specialty. An online, collaborative text may be particularly valuable in this second role. All of the areas, disciplines, roles, and uses of this technology are changing. Indeed, the rate of change seems to increase geometrically, while the number of areas or specialties increases exponentially. Unlike a paper text, an online one can be maintained rapidly and constantly. Once the reader is familiar with the format and content of the work, future reference is easier. So when presented with a new feature, function, or technology there is an updated textbook that's just like the one they used before. While some may prefer a paper copy for initial use, an updated copy may be continually available. Non-Professional Readers. Very few professions today escape being affected by the application of computers and the related communication technology. For those who work in other disciplines an up to date survey and reference may also be valuable when they must deal with the impact of the technology within their own field. When forced to interact with the high priests of technology, they need not compete to get their viewpoint recognized while ignorant of the goals, practices, and terminology of the ordained. When a piece of these technologies interacts with your own business, this will not replace a user manual, or even the level of the "Topic X for Dummies". What it does is to provide the background and context that make those materials more useful. Terminology. There are four words used in this work that require some explanation. The words "platform, entity, element," and "component" are used loosely, but with some intent. Generally the term entity is used to describe the thing being considered or addressed by a topic, element is used to describe a feature or function of the entity being discussed, and component is used to describe a unit that is considered as a whole, and stands alone for some functions, while platform refers to the group of components that are prerequisite to the entity under discussion. Each term frequently represents the specific being discussed and its peers or similar units that operate at the same level. A specific area or topic may be treated at all four levels. Consider a single Web Page as an example. When describing the Client–Server model of the Internet, it is a "component". When considering HTML it is a part of the "platform". In a part of the discussion of Web Site Construction it is first an "element" of the site, and then an "entity" since there are several topics that directly discuss Web Pages. Structure. The question of how to organize a comprehensive technology overview is both fundamental and complex. You can probably think of ways to improve the organization chosen (I know I do every time I look at it). This section outlines of the survey's major organization into parts and breiefly identifies the main subject for each part. The overview is organized into a decomposition hierarchy. The hierarchy is presented two layers at a time. Subject definitions are always presented twice, once with their parent subject and again in their own introduction. "For full detail, see the:" Subject areas and major parts are listed below: Some things fell off this hierarchy, and are contained in other areas of topics. These are: Authors. Any work with the broad and inclusive scope proposed cannot succeed without the contributions of many individuals, both experts and generalists. If this text is to be useful to the general audience, it also needs a lot of amateurs to insure that it stays readable. Thanks in advance to all contributors. As in any good Wiki, details of contribution will be kept for each subject on the page history. Major contributors who wish to be identified will be acknowledged on the CACS/Authors page. There is to be a separate page about writing sections of this book, and it is kept on the CACS/Author Guidelines page. To comment on direction or structure please use the Talk page of the part, chapter, section or topic. 

This Wikibook is a the result of a collaborative effort by a number or contributors. All of them are gratefully acknowledged, and thanked for their good work. The only ones identified here are those who have chosen to share the blame for its shortcomings. You can contact any of them through their Wikibook user talk page or by eMail. Contributors in alphabetic order are: 

Table of Contents. Appendices.  __NOEDITSECTION__ 

Glossary of Computer and Communications Terms For each term in the Glossary, there is a single data page. When a term is usually known by its acronym, it is defined on a page with that title. When needed, redirect pages will be used to get to that data page from any other frequently used names. Many terms contained here have more information than is displayed in a simple definition. When a word or phrase within a definition is shown in "italics", that word or phrase is also expected to appear in the glossary. A list of All terms is contained at the bottom of this page. Terms in the Glossary include: processes, languages, terms of art, acronyms, hardware devices, software package names, companies and organizations. Every effort is made to note those that are trademarks or proprietary names. All such use is intended to be made under "fair use" conditions and should not reduce the owner's rights or the terms of the GNU/GFDL license applied to this work as a whole. The full text of the free use license that applies here is located at Wikipedia License. "The following is Subject to Change. Currently each term is being set up on a subpage of the Glossary. This is mandated by their intended use in pop-up windows when reading the book." List of Terms. A.  -  -  -  -  -  -  - B.  - C.  - /CSS/ G.  - U.  -  -  - 

Hypertext Markup Language (HTML) A document publishing language that mixes content, formatting rules, and processing "scripts" to create "Web Pages". The HTML source is interpreted and displayed by a "browser" on the end user's computer. HTML is an open specification maintained through a consortium of major vendors and corporations known as the "W3C". The most recent, full version of the "language specification" is HTML 4.01 and dates from December 1999. It can be downloaded from the W3C Web Site. HTML (or .htm) is also used as a file extension for a stored web page that uses the language. 

When we talk of set theory, we generally talk about "collections" of certain mathematical objects. In this sense, a "set" can be likened to a bag, holding a finite (or conceivably infinite) amount of things. Sets can be sets of sets as well (bags with bags in them). However, a set cannot contain duplicates -- a set can contain only one copy of a particular item. When we look at sets of certain types of numbers, for example, the natural numbers, or the rational numbers, for instance, we may want to speak only of these sets. These collections of numbers are, of course, very important, so we write special symbols to signify them. We write sets in curly brackets -- { and }. We write all of the "elements", or what the set contains, in the brackets, separated by commas. We generally denote sets using capital letters. For example, we write the set containing the number 0 and the number 1 as {0,1}. If we wish to give it a name, we can say B={0,1}. Special sets. The aforementioned collections of numbers, the naturals, rationals, etc. are notated as the following: Here we will generally write these in standard face bold instead of the doublestruck bold you see above. So we write here N instead of formula_1 (NB following Wikipedia conventions). Notations. We can write some special relations involving sets using some symbols. Containment relations. To say that an element is in a set, for example, 3 is in the set {1,2,3}, we write: We can also express this relationship in another way: we say that 3 is a "member" of the set {1,2,3}. Also, we can say the set {1,2,3} "contains" 3, but this usage is not recommended as it is also used to refer to subsets (see following). We can say that two sets are equal if they contain exactly the same elements. For example, the sets {2,3,1} and {3,1,2} both contain the numbers 1, 2 and 3. We write: We write the set with "no" elements as formula_11, or {}. Here, we use the notation {} for the "empty set" (NB following Wikipedia conventions). The concept of the subset. A very important concept in set theory and other mathematical areas is the concept of the subset. Say we have two sets A={0,1,2,3,4,5,6,7,8,9}, and B={0,1,2,3,4,5}. Now, B "contains some elements" of A, but not all. We express this relationship between the sets A and B by saying B is a "subset" of A. We write this If B is a subset of A, but A is not a subset of B, B is said to be a "proper subset" of A. We write this Note that if formula_13, then formula_12 Intersections and unions. There are two notable and fundamental special operations on sets, the "intersection" and the "union". These are somewhat analogous to multiplication and addition. Intersection. The intersection of two sets A and B are the elements "common" to both sets. For example, if A={1,3,5,7,9} and B={0,1,3}, their intersection, written formula_16 is the set {1,3}. If the intersection of any two sets are empty, we say these sets are "disjoint". Unions. The union of two sets A and B are the "all" elements in "both" sets. For example if A={1,3,5,7,9} and B={0,2,4,6,8}. We say the union, written formula_17 is the set {0,1,2,3,4,5,6,7,8,9}. Set comprehensions. When we write a set, we can do so by writing all the elements in that set as above. However if we wish to write an "infinite" set, then writing out the elements can be too unwieldy. We can solve this problem by writing sets in "set comprehension" notation. We do this by writing these sets including a "rule" and by a relationship to an "index set", say I. That is; where rule can be something like "x"2, or "x"=3"x". For example, this set forms the set of all even numbers: This set forms the set of all solutions to the general quadratic: Universal sets and complements. Universal sets. When we do work with sets, it is useful to think of a larger set in which to work with. For example, if we are talking about sets {-1,0,1} and {-3,-1,1,3}, we may want to work in Z in this circumstance. When we talk about working in such a larger set, such as Z in that instance, we say that Z is a "universal set", and we take all sets to be subsets of this universal set. We write the universal set to be formula_21, however it may be simpler to denote this as E. Complements. Given a set A in a larger universal set E, we define the complement of A to be all elements in E that are not in A, that is the complement of A is: We write the complement as A' or Ac. In this document we will use A'. Problem set. Based on the above information, write the answers to the following questions Answers. 2. No, the square root of 2 is "irrational", not a rational number&lt;br&gt; 4.1. Yes &lt;br&gt; 4.2. No&lt;br&gt; 6. Yes.&lt;br&gt; 8. 5 elements could be {3,5,7,9,11}.&lt;br&gt; 10. formula_37&lt;br&gt; Further ideas. These mentioned concepts are not the only ones we can give to set theory. Key ideas that are not necessarily given much detail in this elementary course in set theory but later in abstract algebra and other fields, so it is important to take a grasp on these ideas now. These may be skipped. Power set. The "power set", denoted P(S), is the set of "all" subsets of S. NB: The empty set is a subset of all sets. For example, P({0,1})= Cardinality. The "cardinality" of a set, denoted |S| is the amount of elements a set has. So |{a,b,c,d}|=4, and so on. The cardinality of a set need not be finite: some sets have infinite cardinality. The cardinality of the power set. If P(S)=T, then |T|=2|S|. Problem set. Based on the above information, write the answers to the following questions. (Answers follow to even numbered questions) Answers. 2. 23=8&lt;br&gt; 4. Set identities. When we spoke of the two fundamental operators on sets before, that of the "union" and the "intersection", we have a set of rules which we can use to simplify expressions involving sets. For example, given: how can we simplify this? Several of the following set identities are similar to those in standard mathematics This is incomplete and a draft, additional information is to be added 

"This page is not yet a normal Index." "The capabilities of search engines now satisfy some of the needs formerly addressed by a paper index." "That means that as this volume nears completion, the content or use of this page may need redesign." "Currently the page is more of a notepad or guide for an author than it is an aid to the ultimate reader." "For any index entry, there is simply a list of those pages, by full name, that have some bearing on the subject." "Having this cross reference will aid in updating topics that involve several pages as well as reviewing them to ensure consistency." "Please Direct comments and suggestions for the format and content of this page to its Talk Page." 

The Polish language is a member of the Western Slavic group of the Indo-European family of languages. It is easiest to learn if one already knows some other related language. The most closely related are other Western Slavic languages: Czech, Slovak, Silesian, Kashubian and Sorbian. More distant are the Southern and Eastern Slavic languages like Bulgarian, Macedonian, Serbo-Croatian, Slovenian and Russian, Ukrainian, Belorusian, respectively. Polish is spoken by a total of approximately 40 million people, making it the second most widely spoken Slavic language in the world, next to Russian. Speakers of related languages can pick it up with not much effort. Someone who doesn't speak any Slavic language, but speaks some other Indo-European language, may still find many similarities between Polish grammar and the grammar of that language, as well as many similar words. This Wikibook is designed for anyone who wants to learn the basics of the Polish language. It is suitable for beginners and those who've been learning the language for a few years. Copyright. This document was originally copyright 2002 Tomasz Węgrzanowski &amp; Anna (grammar) &lt;taw@users.sf.net&gt; It may be distributed under terms of GNU Free Documentation Licence. Since its original writing, it has been edited and redistributed on Wiktionary and, currently, Wikibooks. Further reading.  __NOEDITSECTION__ 

Charles Darwin (1809 - 1882), British naturalist, founder of the Theory of Evolution by means of natural selection. "Note: Darwin's books should be read with a historical mindset, as they do not always reflect current scientific opinion (though impressively, they often do). Scientifically speaking, the knowledge held in his books should be supplemented with other books in . Links. Texts Online: 

Actual Cost of Work Performed (ACWP) The total direct and indirect costs of accomplishing work completed to date. It usually refers either to a given time period or to a completed project phase, activity, or task. Some management techniques compare it to a budget to determine a "variance" or status. In other cases, when it is usually called "earned value", ACWP may be used to determine interim payments to a contractor. 

General Routing Encapsulation (GRE) A method or technique of adding an IP standard header and trailer to a message that does not follow "IP protocols". The encapsulated message is sent over a public network while received messages are stripped of the "wrapper" and processed. This permits non-standard data (from an application like "Notes" or "AppleTalk") and totally encrypted messages to use the Internet. The technology is an important element in "Virtual Private Network" (VPN) offerings. 

Multiprotocol Label Switching (MPLS) is a system of protocols that uses abbreviated routing information based on the devices and connections of a single "Internet Service Provider (ISP)". These abbreviated codes are then attached to an IP message to improve the speed and efficiency of internally routed messages. The protocols and their overall use are documented by the "IETF" in RFC3031. 

Public Key Infrastructure (PKI) An evolving system of "protocols" that define the "Digital Certificates" and "Certificate Authorities" that identify and validate the parties in an "electronic commerce" transaction to each other. 

Tunneling – a software process that uses "General Routing Encapsulation (GRE)" programs to send encrypted or incompatible messages over a network, usually the "Internet". The concept refers to the use of the internet by two nodes of a "Virtual Private Network (VPN)", and their treatment of the internet as a 'tunnel' for their messages. 

Virtual Private Network (VPN) a service, usually offered by an "ISP", that uses various routing, encryption, and security technologies to allow the Internet to serve as the equivalent of a private or controlled network. This supports the connection of a corporate network with external elements for telecommuting or extended enterprise applications, or may entirely replace a proprietary network. VPN is also used to describe the collective software and methods used in the VPN service. 

Adobe Systems Inc. A software vendor specializing in documents. Founded in 1982, their best known products are "Pagemaker, Acrobat Reader," and "Photoshop". Located in San Jose, California, their web address is http://www.adobe.com/main.html. 

Automatic Data Processing (ADP) a corporation, specializing in payroll and "human resources" business services and founded in 1949. As of 2002 they grossed over $15 billion, Headquartered in Roseland, NJ, their web address is www.adp.com.  .ADP is also a file extension for an Access Database Project. 

Introduction. This article examines the concepts of a "function" and a "relation". A relation is any association or link between elements of one set, called the domain or (less formally) the "set of inputs", and another set, called the range or "set of outputs". Some people mistakenly refer to the range as the "codomain"(range), but as we will see, that really means the "set of all possible outputs"—even values that the relation does not actually use. (Beware: some authors do not use the term "codomain"(range), and use the term "range" instead for this purpose. Those authors use the term image for what we are calling "range". So while it is a mistake to refer to the "range" or "image" as the "codomain"(range), it is not necessarily a mistake to refer to "codomain" as "range".) For example, if the "domain" is a set Fruits = {apples, oranges, bananas} and the "codomain"(range) is a set Flavors = {sweetness, tartness, bitterness}, the flavors of these fruits form a relation: we might say that apples are related to (or associated with) both sweetness and tartness, while oranges are related to tartness only and bananas to sweetness only. (We might disagree somewhat, but that is irrelevant to the topic of this book.) Notice that "bitterness", although it is one of the possible Flavors (codomain)(range), is not really used for any of these relationships; so it is not part of the "range" (or "image") {sweetness, tartness}. Another way of looking at this is to say that a relation is a "subset of ordered pairs" drawn from the "set of all possible ordered pairs" (of elements of two other sets, which we normally refer to as the "Cartesian product" of those sets). Formally, R is a relation if for the domain X and codomain(range) Y. The inverse relation of R, which is written as R-1, is what we get when we interchange the X and Y values: Using the example above, we can write the relation in set notation: {(apples, sweetness), (apples, tartness), (oranges, tartness), (bananas, sweetness)}. The inverse relation, which we could describe as "fruits of a given flavor", is {(sweetness, apples), (sweetness, bananas), (tartness, apples), (tartness, oranges)}. (Here, as elsewhere, the order of elements in a set has no significance.) One important kind of relation is the "function". A function is a relation that has exactly one output for every possible input in the domain. (The domain does not necessarily have to include all possible objects of a given type. In fact, we sometimes intentionally use a "restricted domain" in order to satisfy some desirable property.) The relations discussed above (flavors of fruits and fruits of a given flavor) are not functions: the first has two possible outputs for the input "apples" (sweetness and tartness); and the second has two outputs for both "sweetness" (apples and bananas) and "tartness" (apples and oranges). The main reason for not allowing multiple outputs with the same input is that it lets us apply the same function to different forms of the same thing without changing their equivalence. That is, if f is a function with a (or b) in its domain, then a = b implies that f(a) = f(b). For example, z - 3 = 5 implies that z = 8 because f(x) = x + 3 is a function unambiguously defined for all numbers x. The converse, that f(a) = f(b) implies a = b, is not always true. When it is, there is never more than one "input" x for a certain "output" y = f(x). This is the same as the definition of function, but with the roles of X and Y interchanged; so it means the "inverse relation" f-1 must also be a function. In general—regardless of whether or not the original relation was a function—the inverse relation will "sometimes" be a function, and sometimes not. When f and f-1 are both functions, they are called one-to-one, injective, or invertible functions. This is one of two very important properties a function f might (or might not) have; the other property is called onto or surjective, which means, for any y ∈ Y (in the codomain), there is some x ∈ X (in the domain) such that f(x) = y. In other words, a "surjective" function f "maps onto" every possible output at least once. A function can be neither one-to-one nor onto, both one-to-one and onto (in which case it is also called bijective or a one-to-one correspondence), or just one and not the other. (As an example which is neither, consider f = {(0,2), (1,2)}. It is a function, since there is only one y value for each x value; but there is more than one input x for the output y = 2; and it clearly does not "map onto" all integers.) Relations. In the above section dealing with functions and their properties, we noted the important property that all functions must have, namely that if a function does map a value from its domain to its co-domain, it must map this value to only one value in the co-domain. Writing in set notation, if "a" is some fixed value: The literal reading of this statement is: the "cardinality" (number of elements) of the set of all values f(x), such that x=a for some fixed value a, is an element of the set {0, 1}. In other words, the number of "outputs" that a function f may have at any fixed "input" a is either zero (in which case it is "undefined" at that input) or one (in which case the output is unique). However, when we consider the "relation", we relax this constriction, and so a relation may map one value to more than one other value. In general, a relation is any subset of the Cartesian product of its domain and co-domain. All functions, then, can be considered as relations also. Notations. When we have the property that one value is related to another, we call this relation a "binary relation" and we write it as where R is the relation. For arrow diagrams and set notations, remember for relations we do not have the restriction that functions do and we can draw an arrow to represent the mappings, and for a set diagram, we need only write all the ordered pairs that the relation does take: again, by example is a relation and not a function, since both 1 and 2 are mapped to two values, (1 and -1, and 2 and -2 respectively) example let A=2,3,5;B=4,6,9 then A*B=(2,4),(2,6),(2,9),(3,4),(3,6),(3,9),(5,4),(5,6),(5,9) Define a relation R=(2,4),(2,6),(3,6),(3,9) add functions and problems to one another. Some simple examples. Let us examine some simple relations. Say f is defined by This is a relation (not a function) since we can observe that 1 maps to 2 and 3, for instance. Less-than, "&lt;", is a relation also. Many numbers can be less than some other fixed number, so it cannot be a function. Properties. When we are looking at relations, we can observe some special properties different relations can have. Reflexive. A relation is "reflexive" if, we observe that for all values a: In other words, all values are related to themselves. The relation of equality, "=" is reflexive. Observe that for, say, all numbers a (the domain is R): so "=" is reflexive. In a reflexive relation, we have arrows for all values in the domain pointing back to themselves: Note that ≤ is also reflexive (a ≤ a for any a in R). On the other hand, the relation &lt; is not (a &lt; a is false for any a in R). Symmetric. A relation is "symmetric" if, we observe that for all values of a and b: The relation of equality again is symmetric. If "x"="y", we can also write that "y"="x" also. In a symmetric relation, for each arrow we have also an opposite arrow, i.e. there is either no arrow between "x" and "y", or an arrow points from "x" to "y" and an arrow back from "y" to "x": Neither ≤ nor &lt; is symmetric (2 ≤ 3 and 2 &lt; 3 but neither 3 ≤ 2 nor 3 &lt; 2 is true). Transitive. A relation is "transitive" if for all values "a", "b", "c": The relation "greater-than" "&gt;" is transitive. If "x" &gt; "y", and "y" &gt; "z", then it is true that "x" &gt; "z". This becomes clearer when we write down what is happening into words. "x" is greater than "y" and "y" is greater than "z". So "x" is greater than both "y" and "z". The relation "is-not-equal" "≠" is not transitive. If "x" ≠ "y" and "y" ≠ "z" then we might have "x" = "z" or "x" ≠ "z" (for example 1 ≠ 2 and 2 ≠ 3 and 1 ≠ 3 but 0 ≠ 1 and 1 ≠ 0 and 0 = 0). In the arrow diagram, every arrow between two values "a" and "b", and "b" and "c", has an arrow going straight from "a" to "c". Antisymmetric. A relation is "antisymmetric" if we observe that for all values "a" and "b": Notice that antisymmetric is not the same as "not symmetric." Take the relation "greater than or equal to", "≥" If "x" ≥ "y", and "y" ≥ x, then "y" must be equal to "x". a relation is anti-symmetric if and only if a∈A, (a,a)∈R Trichotomy. A relation satisfies "trichotomy" if we observe that for all values "a" and "b" it holds true that: "a"R"b" "or" "b"R"a" The relation "is-greater-or-equal" satisfies since, given 2 real numbers "a" and "b", it is true that whether "a" ≥ "b" or "b" ≥ "a" (both if "a" = "b"). Problem set. Given the above information, determine which relations are reflexive, transitive, symmetric, or antisymmetric on the following - there may be more than one characteristic. (Answers follow.) "x" R "y" if Equivalence relations. We have seen that certain common relations such as "=", and congruence (which we will deal with in the next section) obey some of these rules above. The relations we will deal with are very important in discrete mathematics, and are known as "equivalence relations". They essentially assert some kind of equality notion, or "equivalence", hence the name. Characteristics of equivalence relations. For a relation R to be an "equivalence relation", it must have the following properties, viz. R must be: In the previous problem set you have shown equality, "=", to be reflexive, symmetric, and transitive. So "=" is an equivalence relation. We denote an equivalence relation, in general, by formula_3. Example proof. Say we are asked to prove that "=" is an equivalence relation. We then proceed to prove each property above in turn (Often, the proof of transitivity is the hardest). Thus = is an equivalence relation. Partitions and equivalence classes. It is true that when we are dealing with relations, we may find that many values are related to one fixed value. For example, when we look at the quality of "congruence", which is that given some number "a", a number congruent to "a" is one that has the same remainder or "modulus" when divided by some number "n", as "a", which we write and is the same as writing For example, 2 ≡ 0 (mod 2), since the remainder on dividing 2 by 2 is in fact 0, as is the remainder on dividing 0 by 2. We can show that congruence is an equivalence relation (This is left as an exercise, below Hint use the equivalent form of congruence as described above). However, what is more interesting is that we can group all numbers that are equivalent to each other. With the relation congruence "modulo" 2 (which is using n=2, as above), or more formally: we can group all numbers that are equivalent to each other. Observe: This first equation above tells us all the "even" numbers are equivalent to each other under ~, and all the "odd" numbers under ~. We can write this in set notation. However, we have a special notation. We write: and we call these two sets "equivalence classes". All elements in an equivalence class by definition are equivalent to each other, and thus note that we do not need to include [2], since 2 ~ 0. We call the act of doing this 'grouping' with respect to some equivalence relation "partitioning" (or further and explicitly "partitioning a set S into equivalence classes under a relation ~"). Above, we have partitioned Z into equivalence classes [0] and [1], under the relation of congruence modulo 2. Problem set. Given the above, answer the following questions on equivalence relations (Answers follow to even numbered questions) formula_6 Partial orders. We also see that "≥" and "≤" obey some of the rules above. Are these special kinds of relations too, like equivalence relations? Yes, in fact, these relations are specific examples of another special kind of relation which we will describe in this section: the "partial order". As the name suggests, this relation gives some kind of ordering to numbers. Characteristics of partial orders. For a relation R to be a partial order, it must have the following three properties, viz R must be: We denote a partial order, in general, by formula_7. Question: Example proof. Say we are asked to prove that "≤" is a partial order. We then proceed to prove each property above in turn (Often, the proof of transitivity is the hardest). Reflexive. Clearly, it is true that "a" ≤ "a" for all values a. So ≤ is reflexive. Antisymmetric. If "a" ≤ "b", and "b" ≤ "a", then a "must" be equal to "b". So ≤ is antisymmetric Transitive. If "a" ≤ "b" and "b" ≤ "c", this says that "a" is less than "b" and "c". So "a" is less than "c", so "a" ≤ "c", and thus ≤ is transitive. Thus ≤ is a partial order. Problem set. Given the above on partial orders, answer the following questions Answers. 2. Simple proof; Formalization of the proof is an optional exercise. Posets. A partial order imparts some kind of "ordering" amongst elements of a set. For example, we only know that 2 ≥ 1 because of the partial ordering ≥. We call a set A, ordered under a general partial ordering formula_8, a "partially ordered set", or simply just "poset", and write it (A, formula_8). Terminology. There is some specific terminology that will help us understand and visualize the partial orders. When we have a partial order formula_8, such that "a" formula_8 "b", we write formula_19 to say that a formula_8 but "a" ≠ "b". We say in this instance that a "precedes" b, or "a" is a predecessor of "b". If (A, formula_8) is a poset, we say that "a" is an immediate predecessor of "b" (or "a" immediately precedes "b") if there is no "x" in A such that "a" formula_19 "x" formula_19 "b". If we have the same poset, and we also have "a" and "b" in A, then we say "a" and "b" are "comparable" if "a" formula_8 "b" or "b" formula_8 "a". Otherwise they are "incomparable". Hasse diagrams. "Hasse diagrams" are special diagrams that enable us to visualize the structure of a partial ordering. They use some of the concepts in the previous section to draw the diagram. A Hasse diagram of the poset (A, formula_8) is constructed by Operations on Relations. There are some useful operations one can perform on relations, which allow to express some of the above mentioned properties more briefly. Inversion. Let R be a relation, then its inversion, R-1 is defined by R-1 := {(a,b) | (b,a) in R}. Concatenation. Let R be a relation between the sets A and B, S be a relation between B and C. We can concatenate these relations by defining Diagonal of a Set. Let A be a set, then we define the diagonal (D) of A by Shorter Notations. Using above definitions, one can say (lets assume R is a relation between A and B): R is "transitive" if and only if R • R is a subset of R. R is "reflexive" if and only if D(A) is a subset of R. R is "symmetric" if R-1 is a subset of R. R is "antisymmetric" if and only if the intersection of R and R-1 is D(A). R is "asymmetric" if and only if the intersection of D(A) and R is empty. R is a "function" if and only if R-1 • R is a subset of D(B). In this case it is a function A → B. Let's assume R meets the condition of being a function, then R is "injective" if R • R-1 is a subset of D(A). R is "surjective" if {b | (a,b) in R} = B. Functions. A function is a relationship between two sets of numbers. We may think of this as a "mapping"; a function "maps" a number in one set to a number in another set. Notice that a function maps values to one and only one value. Two values in one set could map to one value, but one value must never map to two values: that would be a relation, "not" a function. For example, if we write (define) a function as: then we say: and we have and so on. This function f maps numbers to their squares. Range and codomain. If D is a set, we can say which forms a née of f is usually a subset of a larger set. This set is known as the "codomain" of a function. For example, with the function f("x")=cos "x", the range of f is [-1,1], but the codomain is the set of real numbers. Notations. When we have a function f, with domain D and range R, we write: If we say that, for instance, "x" is mapped to "x"2, we also can add Notice that we can have a function that maps a point ("x","y") to a real number, or some other function of two variables -- we have a set of ordered pairs as the domain. Recall from set theory that this is defined by the "Cartesian product" - if we wish to represent a set of all real-valued ordered pairs we can take the Cartesian product of the real numbers with itself to obtain When we have a set of "n"-tuples as part of the domain, we say that the function is "n"-ary (for numbers "n"=1,2 we say unary, and binary respectively). Other function notation. Functions can be written as above, but we can also write them in two other ways. One way is to use an arrow diagram to represent the mappings between each element. We write the elements from the domain on one side, and the elements from the range on the other, and we draw arrows to show that an element from the domain is mapped to the range. For example, for the function f("x")="x"3, the arrow diagram for the domain {1,2,3} would be: Another way is to use set notation. If f("x")="y", we can write the function in terms of its mappings. This idea is best to show in an example. Let us take the domain D={1,2,3}, and f("x")="x"2. Then, the range of f will be R={f(1),f(2),f(3)}={1,4,9}. Taking the Cartesian product of D and R we obtain F={(1,1),(2,4),(3,9)}. So using set notation, a function can be expressed as the Cartesian product of its domain and range.  f("x") This function is called "f", and it takes a "variable" "x". We substitute some value for "x" to get the second value, which is what the function maps x to. Types of functions. Functions can either be one to one (injective), onto (surjective), or bijective. "INJECTIVE Functions" are functions in which every element in the domain maps into a unique elements in the codomain. "SURJECTIVE Functions" are functions in which every element in the codomain is mapped by an element in the domain. "'BIJECTIVE Functions" are functions that are both injective and surjective. ---onto functions a function f form A to B is onto , 

Aldus a former software vendor, noted for the developement of "Pagemaker", which was sold to "Adobe" in 1994. 

Animated Cursor – a feature that modifies the displayed "cursor" based on the state of the component the cursor is located in or the properties of a "control" that the cursor points to. State variation usually alters the cursor when a wait state exists. Property variation identifies controls or text that are accessible or available for end-user manipulation.  The "file extension" .ANI is used for the graphics and conditions that control the display. 

Apache A free, open source Web Server "HTTP" processing software package from the "Apache Software Foundation". Apache runs on "OS/2, NetWare, Unix," and "Windows NT". Version 2.0 dates from May 2002. The software source code may be downloaded from the Apache Website, which also contains Apache Documentation. For Apache related questions or issues, you can search this large unofficial Apache Forum hosted by Nabble which currently archives all the Apache projects' mailing lists for cross search and browsing. You can also post your question to the appropriate sub forum which will then forward your post to the corresponding Apache mailing list. 

Apache Software Foundation A collaborative non-profit working group dedicated to the education, development and delivery of open source web software. Its Web address is the Apache Foundation. 

An Application is: Typically, the application uses a single set of libraries and is considered a single project. 

An Assembler is: 

Assembly is the process of creating an executable module from assembly language source code. Sometimes assembly is used as a shorthand reference for "Assembly Language Code". 

Assembly Language is the lowest level of symbolic "progarmming language" characterized by the use of a single "statement" for each machine "instruction". The language also permits the use of labels for routines, variables, and storage assignments. There is a different language for each computer "CPU". Some versions of Assembly Language support the use of "Macro Instructions".  Assembley Language source code is usually stored with a "file extension" of .ASM. 

The Wikibooks Esperanto Textbook is a collaborative project to create an online open-content textbook for , which we hope will become the definitive Esperanto textbook for English speakers. About the book.  __NOEDITSECTION__ 

Wiki Note "This note will be removed later. The entire essay may become a secion within an Introduction Chapter. Or, it might remain here and just have an "assigned reading" on the introduction. Please use the Talk page for this entry for additional discussion. Thanks. =The Myth of the User= There are several terms widely used in the computer industry that are not clearly defined. The result is that everyone understands them and uses them, but the use may tend to complicate rather than illuminate a discussion. This outline looks just two of these terms in the context of an application within a fairly large corporation. Then it draws some conclusions about the impact of context on the use of terminology. Consider the terms User and Application. The User is defined as "the person who uses a computer". The Platform is "the hardware or software that supports an application or a system". The suggested application is one of monitoring and use of a particular resource pool across an area of the company. The Application. Project and unit managers share a common resource pool. They budget and track the use of the pool in their area. They use a series of spreadsheets that show their planned vs. actual use of many resources. The data is extracted for their status reports. The data is also uploaded through the company's networks for use in other summaries of the effective use of the resource pool. These support the manager and planning staff that are responsible for the pool. To make this work, a set of spreadsheet macros are in place to ensure data validity and consistency. These templates are put in place by the corporate IT staff. The "end-user" needs to record their work results, using the spreadsheets by an arbitrary cutoff time on Friday. The Software Hierarchy. Now consider the hierarchy of software engineering that puts this application in place. The example could be extended further. The extensions in Step 3 could be produced by an outside "Value Added Reseller" to produce a resource planning system. Step 1 could look into the earlier history in Hardware for PC Configuration, BIOS, Operating Systems, The C++ compiler, etc. But these 5 layers are quite complex enough to illustrate the terminology problems. "For those who are have never worked in this situation, this is NOT an uncommon example. Millions of people work in this environment every day at your bank, your university, or the manufacturer of your automobile." The Confusion. The above example is fine when it works. But, what about when it fails? As the end-user, I've got a problem in my spreadsheet, and an error message. I find something from the 1st level, the "User's Guide", and try to find a solution. But it is a hopeless quest. I can't use most of the options they discuss. I can't even inspect the VBA code that the guide says my error message comes from. The usually helpful System Admin who put this new update on my desk is equally puzzled, and other users haven't reported a problem. So, what went wrong? Every layer in the hierarchy has their own viewpoint. Everything in the layers above them is viewed as a "Platform" and constitutes a given. Every participant in the layers below them is the User and is considered with the same monolithic view. In some cases, the monolithic view is correct, but even then it is usually distorted. When I buy a new PC for home use, I have the seemingly unfiltered copy of the spreadsheet. But even here, the software was preinstalled by the people I bought the PC from. They chose certain options left by the 1st level spreadsheet team to make the installed software work with this particular configuration. Certain other options may be precluded by their choices. A really good practitioner working in the hierarchy will be partially aware one layer above and below their own level. They can influence the layer above them, through user groups or other methods. They can consider the needs of those one layer below them, and improve their product and documentation to help out, but only as limited by budgets and competitive factors. Conclusions. If you begin to understand the problem, then Welcome to the world of Computers!. But you can also understand how many niches in this world open opportunities for integrators and vendors. How a "Peoplesoft" vendor can package and sell an integrated solution for a business function. How a high-powered consulting firm can charge high prices to help corporate management try to cope. Why the industry keeps looking for something like an "XML" magic wand to make it fit together. And most of all, how some kind of mental framework is needed before you can see how the pieces fit together. This kind of a framework is called an "Architecture", and several of them are needed to see how the pieces fit together. 

Ordinary differential equations involve equations containing: and their solutions. In studying integration, you "already" have considered solutions to very simple differential equations. For example, when you look to solving for g(x), you are really solving the differential equation Notations and terminology. The notations we use for solving differential equations will be crucial in the ease of solubility for these equations. This document will be using three notations primarily: Terminology. Consider the differential equation Since the equation's highest derivative is 2, we say that the differential equation is of "order" 2. Some simple differential equations. A key idea in solving differential equations will be that of integration. Let us consider the second order differential equation (remember that a function acts on a value). How would we go about solving this? It tells us that on differentiating twice, we obtain the constant 2 so, if we integrate twice, we should obtain our result. Integrating once first of all: We have transformed the apparently difficult second order differential equation into a rather simpler one, viz. This equation tells us that if we differentiate a function once, we get formula_9. If we integrate once more, we should find the solution. This is the "solution" to the differential equation. We will get formula_12 for "all" values of formula_13 and formula_14. The values formula_13 and formula_14 are related to quantities known as "initial conditions". Why are initial conditions useful? ODEs (ordinary differential equations) are useful in modeling physical conditions. We may wish to model a certain physical system which is initially at rest (so one initial condition may be zero), or wound up to some point (so an initial condition may be nonzero, say 5 for instance) and we may wish to see how the system reacts under such an initial condition. When we solve a system with given initial conditions, we substitute them after our process of integration. Example. When we solved formula_5 say we had the initial conditions formula_18 and formula_19. (Note, initial conditions need not occur at f(0)). After we integrate we make substitutions: Without initial conditions, the answer we obtain is known as the "general solution" or the solution to the "family of equations". With them, our solution is known as a "specific solution". Basic first order DEs. In this section we will consider "four" main types of differential equations: There are many other forms of differential equation, however, and these will be dealt with in the next section Separable equations. A "separable" equation is in the form (using dy/dx notation which will serve us greatly here) Previously we have only dealt with simple differential equations with g("y")=1. How do we solve such a separable equation as above? We group "x" and "dx" terms together, and "y" and "dy" terms together as well. Integrating both sides with respect to y on the left hand side and x on the right hand side: we will obtain the solution. Worked example. Here is a worked example illustrating the process. We are asked to solve Separating Integrating Letting formula_35 where k is a constant we obtain which is the general solution. Verification. This step does not need to be part of your work, but if you want to check your solution, you can verify your answer by differentiation. We obtained as the solution to Differentiating our solution with respect to x, And since formula_36, we can write We see that we obtain our original differential equation, thus our work must be correct. Homogeneous equations. A "homogeneous" equation is in the form This looks difficult as it stands, however we can utilize the substitution so that we are now dealing with F(v) rather than F(y/x). Now we can express y in terms of v, as "y"="xv" and use the product rule. The equation above then becomes, using the product rule Then which is a separable equation and can be solved as above. However let's look at a worked equation to see how homogeneous equations are solved. Worked example. We have the equation This does not appear to be immediately separable, but let us expand to get Substituting "y"="xv" which is the same as substituting "v"="y"/"x": Now Canceling "v" from both sides Separating Integrating both sides which is our desired solution. Linear equations. A linear first order differential equation is a differential equation in the form Multiplying or dividing this equation by any non-zero function of "x" makes no difference to its solutions so we could always divide by "a"("x") to make the coefficient of the differential 1, but writing the equation in this more general form may offer insights. At first glance, it is not possible to integrate the left hand side, but there is one special case. If "b" happens to be the differential of "a" then we can write and integration is now straightforward. Since we can freely multiply by any function, lets see if we can use this freedom to write the left hand side in this special form. We multiply the entire equation by an arbitrary, "I"("x"), getting then impose the condition If this is satisfied the new left hand side will have the special form. Note that multiplying "I" by any constant will leave this condition still satisfied. Rearranging this condition gives We can integrate this to get We can set the constant "k" to be 1, since this makes no difference. Next we use "I" on the original differential equation, getting Because we've chosen "I" to put the left hand side in the special form we can rewrite this as Integrating both sides and dividing by formula_67 we obtain the final result We call "I" an "integrating factor". Similar techniques can be used on some other calculus problems. Example. Consider First we calculate the integrating factor. Multiplying the equation by this gives or We can now integrate Exact equations. An exact equation is in the form and, has the property that (If the differential equation does not have this property then we can't proceed any further). As a result of this, if we have an exact equation then there exists a function h("x", "y") such that So then the solutions are in the form by using the fact of the total differential. We can find then h("x", "y") by integration Basic second and higher order ODE's. The generic solution of a "n"th order ODE will contain "n" constants of integration. To calculate them we need "n" more equations. Most often, we have either or Reducible ODE's. 1. If the independent variable, "x", does not occur in the differential equation then its order can be lowered by one. This will reduce a second order ODE to first order. Consider the equation: Define Then Substitute these two expression into the equation and we get which is a first order ODE Example. Solve if at "x"=0,  "y"=D"y"=1 First, we make the substitution, getting This is a first order ODE. By rearranging terms we can separate the variables Integrating this gives We know the values of "y" and "u" when "x"=0 so we can find "c" Next, we reverse the substitution and take the square root To find out which sign of the square root to keep, we use the initial condition, D"y"=1 at "x"=0, again, and rule out the negative square root. We now have another separable first order ODE, Its solution is Since "y"=1 when "x"=0, "d"=2/3, and 2. If the dependent variable, "y", does not occur in the differential equation then it may also be reduced to a first order equation. Consider the equation: Define Then Substitute these two expressions into the first equation and we get which is a first order ODE Linear ODEs. An ODE of the form is called linear. Such equations are much simpler to solve than typical non-linear ODEs. Though only a few special cases can be solved exactly in terms of elementary functions, there is much that can be said about the solution of a generic linear ODE. A full account would be beyond the scope of this book If "F(x)=0" for all "x" the ODE is called homogeneous Two useful properties of generic linear equations are Variation of constants. Suppose we have a linear ODE, and we know one solution, "y=w(x)" The other solutions can always be written as "y=wz". This substitution in the ODE will give us terms involving every differential of "z" upto the "n"th, no higher, so we'll end up with an "n"th order linear ODE for "z". We know that "z" is constant is one solution, so the ODE for "z" must not contain a "z" term, which means it will effectively be an "n-1"th order linear ODE. We will have reduced the order by one. Lets see how this works in practice. Example. Consider One solution of this is "y=x2", so substitute "y=zx2" into this equation. Rearrange and simplify. This is first order for D"z". We can solve it to get Since the equation is linear we can add this to any multiple of the other solution to get the general solution, Linear homogeneous ODE's with constant coefficients. Suppose we have a ODE we can take an inspired guess at a solution (motivate this) For this function Dn"y"=pny so the ODE becomes "y=0" is a trivial solution of the ODE so we can discard it. We are then left with the equation This is called the "characteristic" equation of the ODE. It can have up to "n" roots, p1, p2 … pn, each root giving us a different solution of the ODE. Because the ODE is linear, we can add all those solution together in any linear combination to get a general solution To see how this works in practice we will look at the second order case. Solving equations like this of higher order uses exactly the same principles; only the algebra is more complex. Second order. If the ODE is second order, then the characteristic equation is a quadratic, with roots What these roots are like depends on the sign of "b"2-4"c", so we have three cases to consider. "1) b2 &gt; 4c" In this case we have two different real roots, so we can write down the solution straight away. "2) b2 &lt; 4c" In this case, both roots are imaginary. We could just put them directly in the formula, but if we are interested in real solutions it is more useful to write them another way. Defining k2=4c-b2, then the solution is For this to be real, the "A"'s must be complex conjugates Make this substitution and we can write, If "b" is positive, this is a damped oscillation. "3) b2 = 4c" In this case the characteristic equation only gives us one root, "p=-b/2". We must use another method to find the other solution. We'll use the method of variation of constants. The ODE we need to solve is, rewriting "b" and "c" in terms of the root. From the characteristic equation we know one solution is formula_100 so we make the substitution formula_113, giving This simplifies to D2"z"=0, which is easily solved. We get so the second solution is the first multiplied by "x". Higher order linear constant coefficient ODE's behave similarly: an exponential for every real root of the characteristic and a exponent multiplied by a trig factor for every complex conjugate pair, both being multiplied by a polynomial if the root is repeated. E.g., if the characteristic equation factors to the general solution of the ODE will be The most difficult part is finding the roots of the characteristic equation. Linear nonhomogeneous ODEs with constant coefficients. First, let's consider the ODE a nonhomogeneous first order ODE which we know how to solve. Using the integrating factor "e-x" we find This is the sum of a solution of the corresponding homogeneous equation, and a polynomial. Nonhomogeneous ODE's of higher order behave similarly. If we have a single solution, "yp" of the nonhomogeneous ODE, called a "particular" solution, then the general solution is "y=yp+yh", where "yh" is the general solution of the homogeneous ODE. Find "yp" for an arbitrary "F(x)" requires methods beyond the scope of this chapter, but there are some special cases where finding "yp" is straightforward. Remember that in the first order problem "yp" for a polynomial "F(x)" was itself a polynomial of the same order. We can extend this to higher orders. "Example:" Consider a particular solution Substitute for "y" and collect coefficients So "b2=0", "b1=-7", "b0=1", and the general solution is This works because all the derivatives of a polynomial are themselves polynomials. Two other special cases are where "Pn","Qn","An", and "Bn" are all polynomials of degree "n". Making these substitutions will give a set of simultaneous linear equations for the coefficients of the polynomials. 

Basic Language is a high level computer language that was intrroducd with the PC in the 1960s. BASIC is an acronym for Beginner's All-purpose Symbolic Instruction Code. Basic was originally an , but compilers were soon introduced. While there is an "ANSI Standard" Basic, most vendors use proprietary extensions. There is a subset or similar language used in "Web pages" called "Basic Script." Microsoft platform applications include another subset called "Visual Basic for Applications (VBA)". The .BAS "file extension" is frequently used for Basic source code. 

Binary is a number system in the base 2, which contains only the digits 0 and 1. Since computer memory is made up of off-on switches, all computer programs and data are ultimately encoded in binary. Binary refers to a program that has been compiled into an executable module. When the module is stored, the .BIN "file extension" is used. Some development platform compilers produce a binary module that must still be "linked" with other modules to produce an executable program. 

Borland is a software developer and vendor specializing in development tools, language compilers, and data base platforms. It was founded in 1983, and is located in Scotts Valley, California. Their Web address is Borland.com. 

A Computer is programmable machine containing memory, input and output devices, and a "central processor (CPU)"; and capable of executing a stored program. Each computer has a well defined "instruction set" that is used in "programming" it. 

Windows. Language Settings (fastest). The idea of changing the language settings is that you can then type characters quickly and easily (for example by pressing Alt+a for typing á). Windows XP. This is a list that describes how to change the language settings for Windows XP. Now you should see a keyboard icon at your task bar at the bottom. Click on this icon to switch to the United States International keyboard layout. This keyboard layout has a new key (AltGr) and 5 "dead keys". The dead keys are explained below. An interactive diagram of this layout can be found in . The US International keyboard has two different Alt keys. The left Alt key continues to be the regular Alt key, normally associated with Windows menus. The right Alt key becomes what is called AltGr (or graphic Alt) key. This key lets you type very quickly special characters in Spanish and other languages by using "AltGr" and then typing a character from the list below. Note: + indicate typing one key after the other. - is typing two keys at the same time. Lower-Case Characters.  á -&gt; AltGr+a  é -&gt; AltGr+e  í -&gt; AltGr+i  ó -&gt; AltGr+o  ú -&gt; AltGr+u  ü -&gt; AltGr+y  ñ -&gt; AltGr+n  ç -&gt; AltGr+,  å -&gt; AltGr+w Upper-Case Characters.  Á -&gt; AltGr+Shift-a  É -&gt; AltGr+Shift-e  Í -&gt; AltGr+Shift-i  Ó -&gt; AltGr+Shift-o  Ú -&gt; AltGr+Shift-u  Ü -&gt; AltGr+Shift-y  Ñ -&gt; AltGr+Shift-n  Ç -&gt; AltGr+Shift-,  Ð -&gt; AltGr+Shift-d Other Symbols.  ¿ -&gt; AltGr+?  ¡ -&gt; AltGr+!  « -&gt; AltGr+[ Spaniards prefer «angular» quotes  » -&gt; AltGr+] Latin Americans prefer the “curly” ones  ° -&gt; AltGr+: Degree sign; ordinal sign, as in "4° año = cuarto año"  € -&gt; AltGr+5  ¢ -&gt; AltGr+Shift-c  £ -&gt; AltGr+$  ¥ -&gt; AltGr+- Dead Keys. The US International keyboard has five "dead keys". They add the symbol they have marked at the top of the following letter.  '+a = á; '+e = é; ...  "+a = ä; ...; "+u = ü  ~+a = ã; ~+n = ñ  ^+a = â; ...  `+a = à; ...  '+Shift-a = Á; '+Shift+e = É; ...  "+Shift-u = Ü; ...  ~+Shift-n = Ñ; ... To enter what's written on a dead key you need to add a space. ' followed by space generates an actual apostrophe. Other Operating Systems. Information on some other operation systems can be found here. Alt Number Codes. You can type a special character by pressing and holding down the "Alt" button and then typing a number code on the number pad of your keyboard. The most frequently used characters have both a three-digit and a four-digit code. Less frequent characters (such as Á) have only a four-digit code. This page contains a good overview of special characters for different languages Lower Case Characters.  á -&gt; Alt-160 or Alt-0225  ç -&gt; Alt-135 or Alt-0231  é -&gt; Alt-130 or Alt-0233  í -&gt; Alt-161 or Alt-0237  ñ -&gt; Alt-164 or Alt-0241  ó -&gt; Alt-162 or Alt-0243  ú -&gt; Alt-163 or Alt-0250  ü -&gt; Alt-129 or Alt-0252 Upper Case Characters.  Á -&gt; Alt-0193  Ç -&gt; Alt-128 or Alt-0199  É -&gt; Alt-144 or Alt-0201  Í -&gt; Alt-0205  Ñ -&gt; Alt-165 or Alt-0209  Ó -&gt; Alt-0211  Ú -&gt; Alt-0218  Ü -&gt; Alt-154 or Alt-0220 Punctuation Marks.  ¿ Alt-168 or Alt-0191  ¡ Alt-173 or Alt-0161 Copy &amp; Paste. This method can be useful if you are just writing a short text (for example an e-mail) and don't have a computer where you can/want change language settings. Just try to pull up a web page or a document that contains the special characters and paste them into your text. For longer texts, however, this can become quite tedious. Search &amp; Replace. If you are working with a text editor you have the option to search for text and replace it with other text. This feature can be used to 'type' special characters. The idea is to "mark" a character for becoming a special character, for example typing "~a" when you mean "á". After you have written your text you replace marked characters (the "~a") with special characters (the "á"). Of course you have to either type in the Alt number code or paste the character, but the point is that you only have to do it "once" for the whole text and not for every single "á" that you want to type. Automated Search &amp; Replace. If you know a programming language that allows string processing you can automate the "Search &amp; Replace" process by a computer program which automatically replaces all your marked characters with the appropriate special characters after you are done with typing your text. Macintosh. Compared to Windows, typing Spanish characters on a Macintosh is relatively easy. So long as you are using a standard American or UK-style QWERTY keyboard, you may just use the following keyboard commands. (Note that you should release the Option (Opt) key before striking the second letter; for example, for á, hold down Option, strike E, release Option, strike A.) One good way to practice typing Spanish characters on a Mac is to use the Key Caps program, which should be in the Utilities folder in the Applications folder. This simple program will show you what characters you can type next if you hold down the Option and/or Shift keys. Lower-Case Characters.  á -&gt; Opt+E, A  é -&gt; Opt+E, E  í -&gt; Opt+E, I  ó -&gt; Opt+E, O  ú -&gt; Opt+E, U  ü -&gt; Opt+U, U  ñ -&gt; Opt+N, N  ç -&gt; Opt+C  å -&gt; Opt+A Upper-Case Characters.  Á -&gt; Opt+E, Shift+A  É -&gt; Opt+E, Shift+E  Í -&gt; Opt+E, Shift+I  Ó -&gt; Opt+E, Shift+O  Ú -&gt; Opt+E, Shift+U  Ü -&gt; Opt+U, Shift+U  Ñ -&gt; Opt+N, Shift+N  Ç -&gt; Opt+Shift+C Other Symbols.  ¿ -&gt; Opt+Shift+/ (forward slash; same key as ?)  ¡ -&gt; Opt+1  « -&gt; Opt+\ (back slash; under the Delete key)  » -&gt; Opt+Shift+\  *Note that Spaniards prefer «angle quotes,» whereas Latin Americans prefer “curly quotes” as in English.  ° -&gt; Opt+Shift+8 Degree sign; ordinal sign, as in "4° año = cuarto año"  € -&gt; Opt+Shift+2  ¢ -&gt; Opt+4  £ -&gt; Opt+3  ¥ -&gt; Opt+Y KDE/GNOME. In KDE and GNOME you can choose the international US keyboard layout. In KDE, go to Regional &amp; Accessibility - Keyboard Layout in the "KDE Control Center". Add the international US keyboard layout to your active layouts. With the flag icon in your taskbar you can now switch between different layouts. In GNOME, find preferences in the system menu, and select "keyboard" from the preferences menu. Add the US English International layout to your layouts. Alt+Alt will cycle through the various layouts you have added. Alternatively, keyboard layout selection and management may be done via the GNOME "keyboard indicator" applet, which may be added to your panel. With the appropriate keyboard layout selected, you can type á by typing ' and then a:  'a á  'e é  'i í  'o ó  'u ú  "u ü  ~n ñ X Windows. In X Windows (the default windowing system for Linux which both KDE and Gnome run on top of) you can create a custom X keyboard layout with the special Spanish letters added; more information is here. 

Leer en línea (read online). Todos los enlaces llevan a páginas en español ("all links go to Spanish language pages"). Software (software). www.spanishconversationclub.com (Spansih as foreign language conversation club in London) 

Here is a little history about the development of this Wikibook, detailing milestones since the beginning of the project. 

Why Do Substances React? Chemical (and thus, biochemical) reactions only occur to a significant extent if they are energetically favorable. If the products are more stable than the reactants, then in general the reaction will, over time, tend to go forward. Ashes are more stable than wood, so once the energy of activation is supplied (e.g., by a match), the wood will burn. There are plenty of exceptions to the rule, of course, but as a rule of thumb it's pretty safe to say that if the products of a reaction represent a more stable state, then that reaction will go in the forward direction. There are two factors that determine whether or not reactions changing reactants into products are considered to be "favorable": these two factors are simply called enthalpy and entropy. Enthalpy. Simply put, enthalpy is the heat content of a substance ("H"). Most people have an intuitive understanding of what heat is... we learn as children not to touch the burners on the stove when they are glowing orange. Enthalpy is not the same as that kind of heat. Enthalpy is the sum of all the internal energy of a substance's matter plus its pressure times its volume. Enthalpy is therefore defined by the following equation: where (all units given in SI) If the enthalpy of the reactants while being converted to products ends up decreasing (Δ"H" &lt; 0), that means that the products have less enthalpy than the reactants and energy is released to the environment. This reaction type is termed "exothermic". In the course of most biochemical processes there is little change in pressure or volume, so the change in enthalpy accompanying a reaction generally reflects the change in the internal energy of the system. Thus, exothermic reactions in biochemistry are processes in which the products are lower in energy than the starting materials. As an example, consider the reaction of glucose with oxygen to give carbon dioxide and water. Strong bonds form in the products, reducing the internal energy of the system relative to the reactants. This is a highly exothermic reaction, releasing 2805 kJ of energy per mole of glucose that burns (Δ"H" = -2805 kJ/mol). That energy is given off as heat. Entropy. Entropy (symbol "S") is the measure of randomness in something. It represents the most likely of statistical possibilities of a system, so the concept has extremely broad applications. In chemistry of all types, entropy is generally considered important in determining whether or not a reaction goes forward based on the principle that a less-ordered system is more statistically probable than a more-ordered system. What does that mean, really? Well, if the volcano Mt. Vesuvius erupted next to a Roman-Empire era Mediterranean city, would the volcano be more likely to destroy the city, or build a couple of skyscrapers there? It's pretty obvious what would happen (or, rather, what "did" happen) because it makes sense to us that natural occurrences favor randomness (destruction) over order (construction, or in this case, skyscrapers). Entropy is just a mathematical way of expressing these essential differences. When it comes to chemistry, there are three major concepts based on the concept of entropy: Changes in entropy are denoted as ΔS. For the reasons stated above (in the volcano situation), the increase of entropy (ΔS &gt; 0) is considered to be favorable as far as the Universe in general is concerned. A decrease in entropy is generally not considered favorable unless an energetic component in the reaction system can make up for the decrease in entropy (see free energy below). Gibbs Free Energy. Changes of both enthalpy (ΔH) and entropy (ΔS) combined decide how favorable a reaction is. For instance, burning a piece of wood releases energy ("exothermic", favorable) and results in a substance with less structure (CO2 and H2O gas, both of which are less 'ordered' than solid wood). Thus, one could predict that once a piece of wood was set on fire, it would continue to burn until it was gone. The fact that it does so is ascribed to the change in its Gibbs Free Energy. The overall favorability of a reaction was first described by the prominent chemist Josiah Willard Gibbs, who defined the "free energy" of a reaction as where T is the temperature on the Kelvin temperature scale. The formula above assumes that pressure and temperature are constant during the reaction, which is almost always the case for biochemical reactions, and so this book makes the same assumption throughout. The unit of ΔG (for "Gibbs") is the "joule" in SI systems, but the unit of "calorie" is also often used because of its convenient relation to the properties of water. This book will use both terms as convenient, but the preference should really be for the SI notation. What Does ΔG Really Mean? If ΔG &lt; 0 then the reactants should convert into products (signifying a forward reaction)... eventually. (Gibbs free energy says nothing about a reaction's "rate", only its "probability".) Likewise, for a given reaction if ΔG &gt; 0 then it is known that the reverse reaction is favored to take place. A state where ΔG = 0 is called equilibrium, and this is the state where the reaction in both the forward and reverse directions take place at the same rate, thus not changing the net effect on the system. How is equilibrium best explained? Alright, as an example set yourself on the living room carpet with your most gullible younger relative (a little nephew, niece or cousin will work fine). Take out a set of Monopoly, take one ten dollar bill for yourself and give your little relative the rest. Now both of you give the other 5% of all that you have. Do this again, and again, and again-again-again until eventually... you both have the same amount of money. This is precisely what the equilibrium of a reaction means, though equilibrium only very rarely results in an even, 50-50% split of products and reactants. ΔG naturally varies with the concentration of reactants and products. When ΔG reaches 0, the reaction rate in the forward direction and the reaction rate in the reverse direction are the same, and the concentration of reactants and products no longer appears to change; this state is called the "point of chemical equilibrium". You and your gullible little relative have stopped gaining and losing Monopoly money, respectively; you both keep exchanging the same amount each turn. Note again that equilibrium is "dynamic". Chemical reaction does not cease at equilibrium, but products are converted to reactants and reactants are converted to products at exactly the same rate. A small ΔG (that is, a value of ΔG close to 0) indicates that a reaction is somewhat reversible; the reaction can actually run backwards, converting products back to reactants. A very large ΔG (that is, ΔG » 0 or ΔG « 0) is precisely the opposite, because it indicates that a given reaction is irreversible, i.e., once the reactants become products there are very few molecules that go back to reactants. Metabolic pathways. The food we consume is processed to become a part of our cells; DNA, proteins, etc. If the biochemical reactions involved in this process were reversible, we would convert our own DNA back to food molecules if we stop eating even for a short period of time. To prevent this from happening, our "metabolism" is organized in "metabolic pathways". These pathways are a series of biochemical reactions which are, as a whole, irreversible. The reactions of a pathway occur in a row, with the products of the first reaction being the reactants of the second, and so on: At least one of these reactions has to be irreversible, e.g.: The control of the irreversible steps (e.g., A → B) enables the cell to control the whole pathway and, thus, the amount of reactants used, as well as the amount of products generated. Some metabolic pathways do have a "way back", but it is not the same pathway backwards. Instead, while using the reversible steps of the existing pathway, at least one of the irreversible reactions is bypassed by another (irreversible) one on the way back from E to A: This reaction is itself controlled, letting the cell choose the direction in which the pathway is running. Free energy and equilibrium. For ΔG, the free energy of a reaction, standard conditions were defined: Under these standard conditions, ΔG0' is defined as the standard free energy change. For a reaction the ratio of products to reactants is given by keq' (=keq at pH 7.0): The relationship of ΔG0' and keq' is with In theory, we can now decide if a reaction is favorable (ΔG0' &lt; 0). However, the reaction might need a "catalyst" to occur within a reasonable amount of time. In biochemistry, such a catalyst is called an "enzyme". The purpose of DNA melting or DNA denaturation is emphasizing and demonstrating the life cycles of all organisms and the origin of replication. The origin of replication specific structure varies from species to species. Furthermore, the particular sequence of the origin of replication is in a genome which is the human genes. Nevertheless, DNA replication is also part of origin of replication which examen in the living organism such as prokaryotes and eukaryotes. Thermodynamically, there are two important contributions on the DNA denaturation. One of them is the breaking all of the hydrogen bonds between the bases in the double helix; the other one is to overcome the stacking stability/energy of bases on top of each other. There are several methods to denature DNA; heat is known as the most common one use in laboratory. We just have to heat the sample to reach above its melting point, the unstack ability of DNA can be then monitored. Melting point and denaturation of DNA depend on several factors: the length of DNA, base-composition of DNA, the condition of the DNA and also the composition of buffer. For instance, the longer DNA will contain more H-bonds and more intermolecular forces compared to the shorter one; therefore, denaturations of longer DNA requires more time and more heat. Base-composition of DNA can also play as a key factor because A:T requires two hydrogen bonds and G:C interaction requires three hydrogen bonds. The region of DNA which contains more A:T will melt/denature more rapidly compared to G:C. We can also see how the condition of DNA is important because condition of DNA is related to whether the DNA is relax, supercoiled, linear or heavily nicked. It is important because it allow us to examine how much intermolecular forces existing in the double helix. Finally, condition of buffer is also playing an essential role to study DNA denaturation because it allow us to control the amount of ions present in the solution during the entire process. Biologically, DNA denaturation can happen inside the cell during DNA replication or translation. In both cases, DNA denaturation is an essential step and a beginning to start each of the process. Most of the time, denaturation happened because of binding of protein or enzymes to a specific region of DNA, the binding will likely lead to open or denature of the helix. However, the actual meaning of the DNA melting is the denaturation of DNA which changes the structure of DNA from double stranded into single stranded. The processes of DNA denaturation is unwinding the double stranded deoxyribonucleic acid and breaks it into two single stranded by breaking the hydrogen bonding between the bases. DNA denaturation is also known of DNA annealing because it is reservable . The main steps DNA annealing are double helical will go through the denaturation to become partially denatured DNA then it will separated the strands into two single strand of DNA in random coils. 













../Parts of the cell/  Organelles: (1) nucleolus (2) nucleus (3) ribosome (4) vesicle (5) rough endoplasmic reticulum (ER) (6) Golgi apparatus (7) Cytoskeleton (8) smooth ER (9) mitochondrion (10) vacuole (11) cytoplasm (12) lysosome (13) centrioles (14) vacuole Nucleus. The nucleus contains genetic material or DNA in the form of chromatin, or, during mitosis or late interphase, chromosomes. All transcription and replication of genetic material take place within the nucleus, as does RNA processing. The nucleolus also resides within the nucleus and is responsible for RNA transcription and folding. Translation of RNA transcripts takes place outside of the nucleus. Mitochondria. A mitochondrian is the organelle responsible for a cell's metabolism. It synthesizes ATP through a protein called ATP synthase. Mitochondria have a double membrane. An outer membrane and a folded inner membrane. The internal membrane, called the cristae is invaginated (folded or creased), to maximize surface area enabling it to hold more ATP syntheses. It is called as "the powerhouse of the cell" which is present in the eukaryotic organisms. It has matrix inside the inner membrane. It is in a rod shape structure. Ribosomes. Ribosomes are responsible for protein synthesis. They are comprised of interacting protein and nucleic acid chains. Broadly, ribosomes are comprised of a large and a small subunit. The small subunit functions to attach to the mRNA strand and hold it in place during translation, while the large subunit holds and manufactures the growing polypeptide chain. The large subunit is further subdivided into the A (aminoacyl), P (peptidyl), and E (exit) binding sites. Aminoacyl Binding Site The aminoacyl binding site binds a charged tRNA whose anticodon matches the codon in the A site. Peptidyl Binding Site The peptidyl binding site contains the molecular machinery that transfers the bound polypeptide from the tRNA to the polypeptide chain, and holds the growing chain in place. Exit Site The exit site is the terminal binding site for tRNA, where discharged tRNA's are released from the translation complex. Endoplasmic Reticulum. The Endoplasmic Reticulum (ER) acts as a transport from the nucleus and ribosomes to the Golgi apparatus. There are two types of endoplasmic reticulum: Smooth ER. Smooth ER act as transport for various things, mainly the RNA from the nucleus to the ribosomes (RNA is a small piece of the DNA code specifically designed to tell the ribosomes what to make). Smooth ER appears smooth in texture, hence the name. Smooth ER plays an important role in lipid emulsification and digestion in the cell. Rough ER. Rough ER are "rough" because of the ribosomes embedded in them. The rough ER takes the protein to the Golgi apparatus to be packaged into vacuoles Golgi Complex. The Golgi Complex basically functions as a "packaging center" for the cell, attaching "address labels" (functional groups) to various cell products to direct them to their respective locations, and "packaging" the products into vacuoles to ensure delivery. Anatomically, the Golgi Complex consists of layers of lipid membrane stacked one on top of another, with a cis face and a trans face. As the molecular product being packaged moves through the complex, various enzymes act upon it to induce vacuole formation and functional group attachment. &lt;br&gt; Vacuole. Vacuoles are cellular storage places. Like the cell membrane, they are comprised of a lipid bilayer that functions as a selectively permeable barrier to regulate movement of materials into and out of the compartment. They can serve a variety of purposes, storing food, water, or waste products, or immune functions such as containing dangerous materials or maintaining turgor pressure (in plants). Vacuoles serve very different purposes in plant cells than they do in animal cells. Plant Cells In plants, vacuoles comprise a significant portion of the cell's total volume and often contribute significantly to the function of a differentiated cell. For example, vacuoles in stomata cells contain large numbers of potassium ions, which can be pumped in or out to open or close the stomata. Animal Cells In animal cells, vacuoles serve more subordinate roles, such as assisting in endo- and exocytosis or basic storage of food and waste. Central Vacuole The central vacuole is found only in plant cells. It is filled with water and is pressurised, like a balloon. This forces all the other organelles within the cell out toward the cell wall. This pressure is called "turgor pressure" and is what gives plants their "crisp" and firm structure. Peroxisomes. Peroxisomes perform a variety of metabolic processes and as a by-product, produce hydrogen peroxide. Peroxisomes use peroxase enzyme to break down this hydrogen peroxide into water and oxygen. Lysosomes. Lysosomes are vacuoles containing digestive and destructive membranes. In white blood cells, these are used to kill the bacteria or virus, while in tadpole-tail cells they kill the cell by separating the tail from the main body. They also do much of the cellular digestion involved in apoptosis, the process of programmed cell death. 

Cell Biology | Parts of the cell 

../Parts of the cell/ Chloroplasts are the organelles used for photosynthesis (a process that incorporates light energy into storage as chemical energy) whereas mitochondria used in respiration (a process that releases stored chemical energy). It assumed that you already know the information about these organelles explained in the organelles section. If you have not read the entries on chloroplasts and mitochondria from there yet, please go back and read them now. 

Cells are structural units that make up plants and animals; also, there are many single celled organisms. What all living cells have in common is that they are small 'sacks' composed mostly of water. The 'sacks' are made from a phospholipid bilayer membrane. This membrane is semi-permeable (allowing some things to pass in or out of the cell while blocking others). There exist other methods of transport across this membrane that we will get into later. So what is in a cell? Cells are 90% fluid (called cytoplasm) which consists of free amino acids, proteins, carbohydrates, fats, and numerous other molecules. The cell environment (i.e., the contents of the cytoplasm and the nucleus, as well as the way the DNA is packed) affect gene expression/regulation, and thus are VERY important aspects of inheritance. Below are approximations of other components (each component will be discussed in more detail later): Components of cytoplasm. The following is optional reading, as all cell components will be discussed in subsequent chapters. 

Size of Cells. Although it is generally the case that biological cells are too small to be seen at all without a microscope, there are exceptions as well as considerable range in the sizes of various cell types. Eukaryotic cells are typically 10 times the size of prokaryotic cells (these cell types are discussed in the next Chapter). Plant cells are on average some of the largest cells, probably because in many plant cells the inside is mostly a water filled vacuole. So, you ask, what are the relative sizes of biological molecules and cells? The following are all approximations:  0.1 nm (nanometer) diameter of a hydrogen atom  0.8 nm Amino Acid  2 nm Diameter of a DNA Alpha helix  4 nm Globular Protein  6 nm microfilaments  7 nm thickness cell membranes  20 nm Ribosome  25 nm Microtubule  30 nm Small virus (Picornaviruses)  30 nm Rhinoviruses  50 nm Nuclear pore  100 nm HIV  120 nm Large virus (Orthomyxoviruses, includes influenza virus)  150-250 nm Very large virus (Rhabdoviruses, Paramyxoviruses)  150-250 nm small bacteria such as Mycoplasma  200 nm Centriole  200 nm (200 to 500 nm) Lysosomes  200 nm (200 to 500 nm) Peroxisomes  800 nm giant virus Mimivirus  1 µm (micrometer)  (1 - 10 µm) the general sizes for Prokaryotes  1 µm Diameter of human nerve cell process  2 µm E.coli - a bacterium  3 µm Mitochondrion  5 µm length of chloroplast  6 µm (3 - 10 micrometers) the Nucleus  9 µm Human red blood cell  10 µm  (10 - 30 µm) Most Eukaryotic animal cells  (10 - 100 µm) Most Eukaryotic plant cells  90 µm small Amoeba  120 µm Human Egg  up to 160 µm Megakaryocyte  up to 500 µm giant bacterium Thiomargarita  up to 800 µm large Amoeba  1 mm (1 millimeter, 1/10th cm)  1 mm Diameter of the squid giant nerve cell  up to 40mm Diameter of giant amoeba Gromia Sphaerica  120 mm Diameter of an ostrich egg (a dinosaur egg was much larger)  3 meters Length of a nerve cell of giraffe's neck What limits cell sizes? 

The various elements that make up the cell are: The difference between these elements is their respective atomic weights, electrons, and in general their chemical properties. A given element can only have so many other atoms attached. For instance carbon (C) has 4 electrons in its outer shell and thus can only bind to 4 atoms; Hydrogen only has 1 electron and thus can only bind to one other atom. An example would be Methane which is CH4. Oxygen only has 2 free electrons, and will sometimes form a double bond with a single atom, which is an 'ester' in organic chemistry (and is typically scented). As for the organic molecules that make up a typical cell: Here is a list of Elements, symbols, weights and biological roles. 

The question, "What is life?" has been one of many long discussions and the answer may depend upon your initial definitions. Life is cells. Cell theory consists of three basic points. Some definitions of life are: Seven Criteria. In biology, whether life is present is determined based on the following seven criteria: Another way of remembering the seven life processes for children is:-  Movement  Respiration  Sensitivity  Growth  Reproduction  Excretion  Nutrition Note the beginning letter of all the seven life processes, it spells out MRS GREN. Virus Controversy. This definition of life has got some problems to it though: As an example, let's take viruses. Just by your intuition, what would you say: Are viruses alive or dead? Most people's intuitive answer is: Viruses are alive. When we suffer from any viral infection, we have the feeling that these viruses that cause our infection are alive. According to the seven principles as shown above, viruses are dead, as dead as a piece of plastic: They can't reproduce themselves. To understand that, we want to make a quick excursion to the replication mechanism of viruses: Viruses are really strange in their reproduction technique. Humans and other animals reproduce by the means of sexual intercourse, bacteria do something called binary fission: They divide. One cell divides itself into two, the two daughter cells divide again an so on. The point here is that both bacteria and animals or humans reproduce actively without any help from outside. Keep this point in mind as we move on to the viruses. Viruses need other cells to reproduce. They "drill" their way into another cell, called the host cell. Here, they release the genetic material they carry and, by a complex mechanism that shouldn't be explained further at this point, force their host cell to produce exact copies of the virus. After some time, the host cell is full of viruses and bursts, releasing the new viruses into the environment. Thus viruses need help to reproduce. They can't reproduce at all without a host cell and therefore do not fulfill the requirement "It should be able to reproduce itself". Looking at the other parts of the definition we find that viruses maintain some degree of homeostasis (1), being able to keep its protenatious and nucleic machinery separated from the outside world. Viruses also show adaptation(5), with their ability to mutate in order to affect new organisms. In addition to the reproduction problem, they also fail to meet the other requirements, showing no cellular organization (2) (or indeed cells at all), metabolism (3), or growth (4). This example is just to illustrate the problems that arise using this definition. Life is not something one can define as any other technical term in science. Life arose from dead matter around 4 billion years ago. When life can arise from dead matter, there can't be a precise border line between these two. The cell is alive, what about parts of it? Organelles are parts of eukaryotic cells (ones having a nucleus). They help the cell carry out its task. But, are they alive? Do they meet 7 criteria? When a cell divides into two, organelles also 'reproduce'. They also age from young to old and then die. Some of them carry out the task of taking food, converting it to nutrients and energy. They can also react to stimuli, and surely they can evolve. Of course one can argue that all the above are coordinated by the nucleus. But it seems there are some signs of life there. Yes, there are! Scientists have proven that some bacteria, in its evolutionary way, had found a home in other cells. They felt comfortable when living there, and gradually, they have become a part of that cell. Chloroplasts, for example, used to be bacteria. At some point in their evolutionary history these cyanobacteria formed a mutual symbiosis with the proto-eukaryote ancestors of algae. Since that time, chloroplasts have been helping plant cells photosynthesize. Another example is mitochondria, organelles that produce energy for eukaryotes. Very likely a parasitic organism originally, the ancestor of the mitochondria we see today colonized the larger proto-eukarotes. It is unknown if the mitochondrial ancestor originally had a metabolic role in its life cycle or if it adapted to the changing conditions after it was engulfed. 

What makes particularly interesting is that there is so much that is not fully understood. A cell is a complex system with thousands of molecular components working together in a coordinated way to produce the phenomenon we call "life". During the 20th century these molecular components were identified (for example, see Human Genome Project), but research continues on the details of cellular processes like the control of cell division and cell differentiation. Disruption of the normal control of cell division can cause abnormal cell behavior such as rapid tumor cell growth. Cells have complex interactions with the surrounding environment. Whether it is the external world of a single celled organism or the other cells of a multicellular organism, a complex web of interactions is present. Study of the mechanisms by which cells respond appropriately to their environments is a major part of cell biology research and often such studies involve what is called signal transduction. For example, a hormone such as insulin interacting with the surface of a cell can result in the altered behavior of hundreds of molecular components inside the cells. This sort of complex and finely tuned cell response to an external signal is required for normal metabolism and to prevent metabolic disorders like Type II diabetes. Most of the cells of a multi-cellular organism have the same genetic material in every cell; yet, there may be hundreds of different types of cells that make up the organism's body each with its own distinctive shape, size, and function. In any case, all of these cells were developed from one special cell, a zygote. The study of how the many cell types develop during embryonic development (Developmental Biology) is a branch of Biology that is heavily dependent on the use of microscopy. Much of the control of cell differentiation is at the level of the control of gene transcription, the control of which mRNAs are made. Muscle cells make muscle proteins and nerve cells make brain proteins. Geneticists, molecular biologists and cell biologists are working to discover the details of how cells specialize to accomplish hundreds of functions from muscle contraction to memory storage. 

 Most of these prokaryotic cells are small, ranging from 1 to 10 microns with a diameter no greater than 1 micron. The major differences between Prokaryotic and Eukaryotic cells are that prokaryotes do not have a nucleus as a distinct organelle and rarely have any membrane bound organelles [mitochondria, chloroplasts, endoplasmic reticulum, golgi apparatus, a cytoskeleton of microtubules and microfilaments] (the only exception may be a bacterium discovered to have vacuoles). Both types contain DNA as genetic material, have a surrounding cell membrane, have ribosomes[70 s], accomplish similar functions, and are very diverse. For instance, there are over 200 types of cells in the human body, that vary greatly in size, shape, and function. Prokaryotes are cells without a distinct nucleus.They have genetic material but that material is not enclosed within a membrane. Prokaryotes include bacteria and cyanophytes. The genetic material is a single circular DNA strand and is located within the cytoplasm. Recombination happens through transfers of plasmids (short circles of DNA that pass from one bacterium to another). Prokaryoytes do not engulf solids, nor do they have centrioles or asters. Prokaryotes have a cell wall made up of peptidoglycan. In majority of prokaryotes, the genome consists of a circular chromosome whose structure includes fewer proteins that found in the linear chromosomes of eukaryotes. Their chromosome is located in the nucleoid, a region of cytoplasm that appears lighter than surrounding cytoplasm in electron micrographs. Also, a single chromosome have much smaller rings of separately replication DNA called plasmids. Cell Surface. Prokaryotic cell walls maintain cell shape, provide physical protection, and prevents the cell from bursting in a hypotonic environment. In hypertonic environment, most prokaryotes lose water and shrink away from their wall (plasmolyze). The cell walls of prokaryotes differ in molecular composition and construction from those of eukaryotes. The bacterial cell walls contain peptidoglycan, a network of modified-sugar polymers cross linked by short polypeptides. This molecular fabric encloses the entire bacterium and anchors other molecules that extend from its surface. Archaeal cell walls contain a variety of polysaccharides and proteins but lack peptidoglycan. Gram-positive bacteria have simpler walls with a relatively large amount of peptidoglycan. It has a thick cell wall that traps the crystal violet in the cytoplasm. The alcohol rinse does not remove the crystal violet which masks the added red safanin dye. Gram-negative bacteria have less peptidoglycan and are structurally more complex, with an outer membrane that contains lipopolysaccharides. It has a thinner layer of peptidoglycan, and it is located in a layer between the plasma membrane and an outer membrane. The crystal violet is easily rinsed from the cytoplasm, and the cell appears pink or red. The cell wall of many prokaryotes is covered by a capsule, a sticky layer of polysaccharide or protein. The capsule enables prokaryotes to adhere to their substrate or to other individuals in a colony. Some capsules protect against dehydration, and some shield pathogenic prokaryotes from attack by their host's immune system. Some prokaryotes stick to their substrate or to one another by means of hair like protein appendages called fimbriae. They are also known as attachment pili. Fimbriae are usually shorter extension of the plasma membrane. In uniform environment, flagellated prokaryotes move randomly, but in heterogeneous environment, many prokaryotes exhibit taxis, movement toward or away from a stimulus. For example, prokaryotes that exhibit chemotaxis change their movement pattern in response to chemicals. They move toward nutrients or oxygen (positive chemotaxis) or away from a toxic substance (negative chemotaxis). Reproduction and Adaptation. Prokaryotes reproduce quickly in a favorable environment. By binary fission, a single prokaryotic cell divid into 2 cells, which then divide into 4, 8, 16, and on. Under optimal conditions, many prokaryotes can divide every 1-3 hours. However the cells eventually exhaust their nutrient supply, poison themselves with metabolic wastes, face competition from other microorganisms, or are consumed by other organisms. The prokaryotes are small, they reproduces by binary fission, and they have short generation times. The ability of some prokaryotes to withstand harsh conditions also contributes to their success. Certain bacteria develop resistant cell called endospores when an essential nutrient is lacking. The original cell produces a copy of its chromosome and surrounds it with a tough wall, forming the endospore. Water is removed from the endospore, and its metabolism halt. The rest of the original cell then disintegrates, leaving the endospore behind. Most endospore are so durable that they can survive in boiling water. In less hostile environments, endospore can remain dormant but viable for centuries, able to rehydrate and resume metabolism when their environment improves. Due to their short generation times, prokaryotic populations can evolve substantially in short periods of time. The ability of prokaryotes to adapt rapidly to new conditions highlights the fact that although the structure of their cells is simpler than that of eukaryotic cells, prokaryotes are not "primitive" or "inferior" in an evolutionary sense. They are highly evolved, and their population have responded successfully to many different types of environmental challenges. Rapid reproduction and mutation In sexually reproducing species, the generation of a novel allele by a new mutation is rare at any particular gene. Instead, most of the genetic variation in sexual populations results from the way existing alleles are arranged in new combinations during meiosis and fertilization. Prokaryotes do not reproduce sexually, so at first glance their extensive genetic variation may seem puzzling. After repeated rounds of division, most of the offspring cells are genetically identical to the original parent cell; however owing to insertions, deletions, and base-pair substitutions in their DNA, some of the offspring cells may differ genetically. The new mutations, though individually rare, can greatly increase genetic diversity in specie that has short generation times and large population sizes. This diversity, in turn, can lead to rapid evolution: individuals that are genetically better equipped for their local environment tend to survive and reproduce more prolifically than less fit individuals. Transformation and Transduction. In transformation, the genotype and possible phenotype of a prokaryotic cell are altered by the uptake of foreign DNA from its surroundings. For example, bacteria from a harmless strain of Streptococcus pneumonia can be transformed to pneumonia-causing cells if they are placed into a medium containing dead, broken-open cells of the pathogenic strain. This transformation occurs when a live nonpathogenic cell takes up a piece of DNA carry the allele for pathogenicity. The foreign allele is then incorporated into the cell's chromosome, replacing the existing nonpathogenic allele- an exchange of homologous DNA segments. The cell is now a recombinant: Its chromosome contains DNA derived from two different cells. In transduction, bacteriophage carries bacterial genes from one hose cell to another; transduction is a type of horizontal gene transfer. For most phages, transduction results from accidents that occur during the phage reproductive cycle. A virus that carries bacterial DNA may not be able to reproduce because it lacks its own genetic material. However, the virus may be able to attach to another bacterium (a recipient) and inject the piece of bacterial DNA acquired from the first cell (the donor). Some of this DNA may subsequently replace the homologous region of the recipient cell's chromosome by DNA recombination. In such a case, the recipient cell's chromosome becomes a combination of NA derived from two cells; genetic recombination has occurred.  Conjugation and Plasmids In a process called conjugation, genetic material is transferred between two bacterial cells ( of same or different species) that are temporarily joined. The DNA transfer is one way: One cell donates the DNA, and the other receives it. The donor uses sex pili to attach to the recipient. After contacting a recipient cell, each sex pilus retracts, pulling the two cells together, much like a grappling hook. A temporary "mating bridge" then forms between the two cells, providing an avenue for DNA transfer. In most cases, the ability to form sex pili and donate DNA during conjugation results from the presence of a particular piece of DNA called the F factor. The F factor consists about 25 genes, most required for the production of sex pili. The F factor can exist either as a plasmid or as a segment of DNA within the bacterial chromosome. The F factor in its plasmid form is called F plasmid. Cells containing the F plasmid, designated F+ cells, function as DNA donors during conjugation. Cells lacking the F factor, designated F-, function as DNA recipients during conjugation. The F+ condition is transferable in the sense that an F+ cell converts and F- cell to F+ is a copy of the entire F+ plasmid is transferred. Chromosomal genes can be transferred during conjugation when the donor cell's F factor is integrated into the chromosome. A cell with the F factor built into its chromosome is called an Hfr cell. Like an F+ cell, an Hfr cell functions as a donor during conjugation with an F- cell. When chromosomal DNA from an Hfr cell enters and F- cell, homologous regions of the HFr and F- chromosomes may align, allowing segments of their DNA to be exchanged. This results in the production of a recombinant bacterium that has genes derived from two different cells- a new genetic cariant on which evolution can act. Though these processes of horizontal gene transfer have so far been studied almost exclusively in bacteria, it is assumed that they are similarly important in archaea. Diverse nutritional and metabolic adaptations. The mechanisms discussed in the previous section- rapid reproduction, mutation, and genetic recombination- underlie that extensive genetic variation found in prokaryotic populations. This variation is reflected in the nutritional adaptations found in prokaryotes. Like all organisms, prokaryotes can be categorized by their nutrition; how they obtain every and the carbon used in building the organic molecules that make up cells. Nutritional diversity is greater among prokaryotes than among eukaryotes: Every type of nutrition observed in eukaryotes is represented among the prokaryotes, along with some nutritional modes unique to prokaryotes. Phototrophs are the organisms that obtain energy from light. Chemotrophs are the organisms that obtain energy from chemicals. Organisms that need only an inorganic compound are called autotrophs. Heterotrophs require at least one organic nutrient to make other organic compounds. Combining these possibilities for energy sources and carbon sources results in four major modes of nutrition.  Photoautotrophs: photosynthetic organisms that capture light energy and use it to drive the synthesis of organic compounds and other inorganic carbon compounds. Cyanobacteria and many other groups of prokaryotes are photoautotrophs, as are plants and algae. Chemoautotrophs: also need only an inorganic compound; however, instead of using light as an energy source, they oxidize inorganic substance, such as hydrogen sulfide, ammonia, or ferrous ions. This mode of nutrition is unique to certain prokaryotes. Photoheterotrophs: Harness energy from light but must obtain carbon in organic form. This mode is unique to certain marine and halophilic (salt-loving) prokaryotes. Chemoheterotrphs: must consume organic molecules to obtain both energy and carbon. This nutritional mode is widespread among prokaryotes. Fungi, animals, most protists, and even some parasitic plants are also chemoheterotrophs. The Role of Oxygen In Metabolism. Prokaryotic metabolism also varies with respect to oxygen. Obligate aerobes use oxygen for cellular respiration and cannot grow without it. Obligate anaerobes, however, are posioned by oxygen. Some obligate anaerobes live exclusively by fermentation; other extract chemical energy by anaerobic respiration, in which substance other than oxygen such as nitrate ions or sulfate ions accept electrons at the "downhill" end of electron transport chains. Facultative anaerobes use oxygen if it is present but can also carry out anaerobic respiration or fermentation in an anaerobic environment.  Nitrogen Metabolism Nitrogen is essential for the production of amino acids and nucleic acids in all organisms. Whereas eukaryotes can obtain nitrogen from only a limited group of nitrogen compounds, prokaryotes can metabolize nitrogen in a wide variety of forms. For example, some cyanobacteria and some methanogens covert atmospheric nitrogen to ammonia, a process called nitrogen fixation. The cells can then incorporate this "fixed" nitrogen into amino acids and other organic molecules. In terms of their nutrition, nitrogen-fixing cyanobacteria are some of the most self-sufficient organisms, since they need only light, carbon dioxide, nitrogen, water and some minerals to grow. Nitrogen fixation by prokaryotes has a large impact on other organisms. For example, nitrogen -fixing prokaryotes can increase the nitrogen but can use the nitrogen compounds that the prokaryotes produce from ammonia.  Metabolic Cooperation Cooperation between prokaryotes allows them to use environmental resource they could not use as individual cells. In some cases, this cooperation takes place between specialized cells of a colony. For instance, the cyanobacterium Anabaena has genes than encode proteins for photosynthesis and for nitrogen fixation, but a single cell cannot carry out both processes at the same time. The reason is that photosynthesis produces oxygen which inactivates the enzymes involved in nitrogen fixation. Instead of living as isolated cells, anabaena forms filamentous colonies synthesis while a few specialized cells called heterocytes carry out only nitrogen fixation. Each heterocyte is surrounded by a thickened cell wall that restricts entry of oxygen produced by neighboring photosynthetic cells. Intercellular connections allow heterocytes to transport fixed nitrogen to neighboring cells and to receive carbohydrates. Metabolic cooperation between different prokaryotic species often occurs in surface-coating colonies known as biofilms. Cells in a biofilm secrete signaling molecules that recruit nearby cells, causing the colonies to grow. The cells also produce proteins that stick the cells to the substrate and on to another. Channels in the biofilm allow nutrients to reach cells in the interior and wastes to be expelled. Biofilm damage industrial and medical equipment, contaminate products, and contribute to tooth decay and more serious health problems. In another example of cooperation between prokaryotes, sulfate-consuming bacteria coexist with methane-consuming archaea in ball- shaped aggregates on the ocean floor. The bacteria appear to use the archaea's waste products, such as organic compounds and hydrogen. In turn, the bacteria produce compounds that facilitate methane consumption by the archaea. This partnership has global ramifications. Reference. Berg, Jeremy M., John L. Tymoczko, and Lubert Stryer. Biochemistry. 7th ed. New York: W.H. Freeman, 2012. Print. Reece, Campbell, Lisa A. Urry, Michael L. Cain, Steven A. Wasserman, Peter V. Minosky, and Robert B. Jackson. Biology. 8th ed. San Francisco: Cummings, 2010. Print. See also. Eukaryotes 

Eukaryotes house a distinct nucleus, a structure in which the genetic material (DNA) is contained, surrounded by a membrane much like the outer cell membrane. Eukaryotic cells are found in most algae, protozoa, all multicellular organisms (plants and animals) including humans. The genetic material in the nucleus forms multiple chromosomes that are linear and complexed with proteins that help the DNA 'pack' and are involved in regulation of gene expression. The cells of higher plants differ from animal cells in that they have large vacuoles, a cell wall, chloroplasts, and a lack of lysosomes, centrioles, pseudopods, and flagella or cilia. Animal cells do not have the chloroplasts, and may or may not have cilia, pseudopods or flagella, depending on the type of cell. Comparing Prokaryotic and Eukaryotic Cells. All cells have several basic features in common: they are all bounded by a selective barrier, plasma membrane. Cytosol is a jellylike substance that is semifluid. All cells contain chromosomes which carry genes in the form of DNA, and ribosomes that make proteins according to instructions from the gene. The major difference between prokaryotic and eukaryotic cells is the location of their DNA. In eukaryotic cell, DNA is found at the nucleus, which is bounded by a double membrane. (the word eukaryotic is from the Greek eu, true, and karyon, kernel, here referring to the nucleus). Eukaryotic cells are much larger than prokaryotic cells; size is general aspect of cell structure that relates to function. The logistics of carrying out cellular metabolism sets limits on cell size. At the lower limit, the smallest cells, known are bacteria called mycoplasmas have diameters between 0.1 and 1.0mm. These are the smallest packages with enough DNA to program metabolism and enough enzymes and other cellular equipment to carry out the activities necessary for a cell to sustain itself and reproduce. Metabolic requirements also impose theoretical upper limits on the size that is practical for a singel cell. Plasma membrane functions as a selective barrier that allows sufficient passage of oxygen, nutrients, and wastes to service the entire cell. For each square micrometer of membrane, only a limited amount of a particular substance can cross per second, so the ratio of surface area to volume is critical. As a cell increases in size, its volume grows proportionately more than its surface area. Area is proportional to a linear dimension squared, whereas volume is proportional to the linear dimension cubed. Therefore a smaller object has a greater ration of surface area to volume. The need for a surface area sufficiently large to accommodate the volume helps explain the microscopic size of most cells, and the narrow, elongated shapes of others, such as nerve cells. Larger organisms has more cells compare to smaller cells. High ratio of surface area to volume is especially important in cells that exchange a lot of material with their surroundings such as intestinal cells. Such cells may have many long, thin projections from their surface called microvilli, which increase surface area without an appreciable increase in volume. Animal Cells. Flagellum: locomotion organelle present in some animal cells; composed of a cluster of microtubules within an extension of the plasma membrane. Centrosome: region where the cell's microtubules are initiated contains a pair of centrioles which function is unknown. Cytoskeleton: reinforces cell's shape, functions in cell movement components are made of protein. It includes microfilaments, intermediate filaments, and microtubules. Microvilli: projections that increase the cell's surface area. Peroxisome: organelle with carious specialized metabolic functions; produces hydrogen peroxide as a by-product, then converts it to water. Mitochondrion: organelle where cellular respiration occurs and most ATP is generated. Lysosome: digestive organelle where macromolecules are hydrolyzed. Golgi apparatus: organelle active in synthesis, modification, sorting, and secretion of cell products. Ribosomes: complexes (small brown dots) that make proteins; free in cytosol or bound to rough ER or nuclear envelope. Plasma membrane: membrane enclosing the cell Endoplasmic Reticulum (ER): network of membraneous sacs and tube; active in membrane synthesis and other synthetic and metabolic processes; has rough (ribosome-studded) and smooth regions. (Rough ER, and Smooth ER) Nucleus: nucleus contains:  "Nuclear envelope": double membrane enclosing the nucleus; perforated by pores; continuous with ER  "Nucleolus": structure involved in production of ribosomes; a nucleus has one or more nucleoli  "Chromatin": material consisting of DNA and proteins; visible as individual chromosomes in a dividing cell In animal cells, lysosomes, centrosomes with centrioles, and flagella are present but not in plant cells. Plant Cell. Cell Wall: outer layer that maintains cell's shape and protects cell from mechanical damage; made of cellulose, other polysaccharide, and protein. Plasmodesmata: channels through cell walls that connect the cytoplasms of adjacent cells. Chloroplast: photosynthetic organelle; converts energy of sunlight to chemical energy stored in sugar molecules.  Central vacuole: prominent organelle in older plant cells; functions include storage, breakdown of waste products, hydrolysis of macromolecules; enlargement of vacuole is a major mechanism of plant growth. Nucleus: nucleus contains:  "Nuclear envelope": double membrane enclosing the nucleus; perforated by pores; continuous with ER  "Nucleolus": structure involved in production of ribosomes; a nucleus has one or more nucleoli  "Chromatin": material consisting of DNA and proteins; visible as individual chromosomes in a dividing cell Golgi apparatus: organelle active in synthesis, modification, sorting, and secretion of cell products. Endoplasmic Reticulum (ER): network of membraneous sacs and tube; active in membrane synthesis and other synthetic and metabolic processes; has rough (ribosome-studded) and smooth regions. (Rough ER, and Smooth ER) Ribosomes: complexes (small brown dots) that make proteins; free in cytosol or bound to rough ER or nuclear envelope. Cytoskeleton: reinforces cell's shape, functions in cell movement components are made of protein. It includes microfilaments, intermediate filaments, and microtubules. In plant cell, chloroplasts, central vacuole, cell wall, and plasmodesmata are present but not in animal cells. Chromatin in the plant cell is a primary protein Nucleus. The nucleus contains most of the genes in the eukaryotic cell; some genes are located in mitochondria and chloroplast. It is generally the most conspicuous organelle in a eukaryotic cell. The nuclear envelope encloses the nucleus, sparating its contents from the cytoplasm. The nuclear envelope is a double membrane, each a lipid bilayer with associated proteins. The envelope is perforated by pore structure that are about 100nm in diameter. At the lip of each pore, the inner and outer membranes of the nuclear envelope are continuous. Pore complex lines each pore and regulates the entry and exit of most proteins and RNAs, as well as large complexes of macromolecules. Except at the pores, the nuclear side of the envelope is lined by the nuclear lamina, a netlike array of protein filaments that maintains the shape of the nucleus by mechanically supporting the nuclear envelope. Also nuclear matrix, a framework of fibers extending throughout the nuclear interior, present. Chromosomes are organized DNA units that carry the genetic information. Each chromosome is made up of material called chromatin, a complex of proteins and DNA. Stained chromatic usually appears as a diffuse mass, byt as a cell prepares to divide, the thin chromatin fibers coil up and condense thick enough to be distinguished as chromosomes. Each eukaryotic species has a characteristic number of chromosomes. For example human has 46 chromosomes. Nucleolus is a prominent structure within the nondividing nucleus. Ribosomal RNA (rRNA) is synthesized from instructions in the DNA; in nucleolus, proteins imported from the cytoplasm are assembled with rRNA into large and small ribosomal subunits. Theses subunits then exit the nucleus through the nuclear pores to the cytoplasm, where a large and a small subunit can assemble into a ribosome. the number depends on the species and the stage in the cell's reproductive cycle. The Nucleus directs protein synthesis by synthesizing messenger RNA (mRNA) according to instructions provided by the DNA. The mRNA is then transported to the cytoplasm via the nuclear pores. Once an mRNA molecule reaches the cytoplasm, ribosomes translate the mRNA's genetic message into the primary structure of a specific poly peptide. Ribosomes. Ribosomes are complexes made of ribosomal RNA and protein; ribosomes are the cellular components that carry out proteins synthesis, also known as protein factories. Cells that have high rates of protein synthesis have particularly large number of ribosomes. Cells active in protein synthesis also have prominent nucleoli. Ribosomes build proteins in two cytoplasmic locales. Free ribosomes are suspended in the cytosol, while bound ribosomes are attached to the outside of the endoplasmic reticulum or nuclear envelope. Bound and free ribosomes are structurally identical, and ribosomes can alternate between the two roles. Most of proteins are made on free ribosomes function within the cytosol. Bound ribosomes generally make proteins that are destined for insertion into membranes, for packaging within certain organelles such as lysosomes, or for export from the cell (secretion). The Endomembrane System. Endomembrane system carries out a variety of tasks in the cell. These tasks include synthesis of proteins and their transport into membranes and organelles or out of the cell, metabolism and movement of lipids, and detoxification of poisons. The membrane of this system are related either through direct physical continuity or by the transfer of membrane segments as tiny vesicles. The various membranes are not identical in structure and function; the thickness, molecular composition, and types of chemical reactions carried out in a given membrane are not fixed but modified several times during the membrane's life. The endomembrane system includes the nuclear envelope, the endoplasmic reticulum, the Golgi apparatus, lysosomes, various kinds of vacuoles, and the plasma membrane. Endoplasmic Reticulum (ER). Endoplasmic reticulum (ER) is an extensive network of membrane that it accounts for more than half the total membrane in many eukaryotic cells. The word endoplasmic means "within the cytoplasm", and reticulum is Latine for "little net". The ER consists of a network of membranous tubules and sacs called cisternae. The ER membrane separates the internal compartment of the ER, ER lumen (cavity) or cisternal space, from the cytosol. Since ER membrane is continuous with the nuclear envelope, the space between the two membranes of the envelope is continuous with the lumen of the ER. Smooth ER lacks ribosomes on its outer surface, and Rough ER has ribosomes on the outer surface of the membrane. Ribosomes are also attached to the cytoplasmic side of the nuclear envelope's outer membrane. Smooth ER- The smooth ER functions in diverse metabolic processes, which vary with cell type. These processes include synthesis of lipids, metabolism of carbohydrates, and detoxification of drugs and poisons. Enzymes of the smooth ER are important in the synthesis of lipids, including oils, phospholipids, and steroids. Sex hormones of vertebrates and the various steroid hormones are produced by the smooth ER in animal cells. Other enzymes of the smooth ER help detoxify drugs and poisons in liver cells. Detoxification involves adding hydroxyl groups to drug molecules, making them more soluble and easier to flush from the body. For example, sedative phenobarbital and other barbiturates are the drugs that metabolized in this manner by smooth ER in liver cells. Barbiturates, alcohol, and many other drugs induce the proliferation of smooth ER and its associated detoxification enzymes, therefore, increasing tolerance to the drugs; in other words, higher doses are required to achieve a particular effect. Also, because some of the detoxification enzymes have relatively broad action, the proliferation of smooth ER in response to one drug can increase tolerance to other drugs as well. The smooth ER also stores calcium ions; in muscle cells, a specialized smooth ER membrane pumps calcium ions from the cytosol into the ER lumen. When a muscle cell is stimulated by a nerve impulse, calcium ions rush back across the ER membrane into the cytosol and trigger contraction of the muscle cell. Rough ER- Many times of cells secrete proteins produced by ribosomes attached to rough ER. As a polypeptide chain grows from a bound ribosomes, it is threaded into the ER lumen through a pore formed by a protein complex in the ER membrane. As the new protein enters the ER lumen, it folds into its native shape. Most secretory proteins are glycoproteins, which have carbohydrates covalently bonded to them. After secretory proteins are formed, the ER membrane keeps them separate from proteins that are produced by free ribosomes and will remain in the cytosol. Secretory proteins depart from the ER wrapped in the membranes of vesicles that bud like bubbles from a specialized region called transitional ER. Transport vesicles are the vesicles in transit from one part of the cell to another. Rough ER is also a membrane factory for the cell; it grows in place by adding membrane proteins and phospholipids to its own membrane. As polypeptide destined to be membrane proteins grow from the ribosomes, they are inserted into the ER membrane and are anchored there by their hydrophobic portions. The rough ER makes its own membrane phospholipids; enzymes build into the ER membrane assemble phospholipids from precursors in the cytosol. The ER membrane expands and is transferred in the form of transport vesicles to other components of the endomembrane system. Golgi Apparatus. Golgi is a center of manufacturing, warehousing, sorting, and shipping. The products of the ER are modified and stored and then sent to other destinations. Golgi apparatus is extensive in cells specialized for secretion. The Golgi apparatus consists of flattened membranous sac, cisternae. The membrane of each cisterna in a stack separates ints internal space from the cytosol. Besicles concentrated in the vicinity of the Golgi apparatus are engaged in the transfer of material between parts of the Golgi and other structures. Golgi stack has a distinct structural polarity with the membrane of cisternae on opposite side of the stack different in thickness and molecular composition. The two poles of a Golgi stack are referred to as the cis face and the trans face; cis is the receiving and trans is shipping departments of the Golgi apparatus. The cis face is usually located near ER. Transport vesicles move material from the ER to the Golgi apparatus. A vesicle that buds from the ER can add its membrane and the contents of its lumen to the cis face by fusing with a Golgi membrane. The trans face give rise to vesicles, which pinch off and travel to other sites. The products of ER are usually modified during their transit from the cis region to the trans region of the Golgi. Various Golgi enzymes modify the carbohydrate portions of glycoproteins; carbohydrates are first added to proteins in the rough ER during the process of polypeptide synthesis. The carbohydrate on the resulting glycoprotein is then modified as it passes through the rest of the ER and the Golgi. The Golgi removes some sugar monomers and substitutes other, producing a large variety of carbohydrates. In addition, the Golgi apparatus manufactures certain macromolecules by itself. Many polysaccharides secreted by cells are Golgi products, including pectins and certain other non-cellulose polysaccharides made by plant cells and incorporated along with cellulose into their cell walls. Similar to secretory proteins, non-protein Golgi products will be secreted depart from the trans face of the Golgi inside transport vesicles that eventually fuse with the plasma membrane. The Golgi manufactures and refines its products in stages, with different cisternae containing unique teams of enzymes. Recent research has give rise to a new model of the Golgi as a more dynamic structure; According to the cisternal maturation model, the cisternae of the Golgi actually progress forward from the cis to the tras face of the Golgi, carrying and modifying their cargo as they move. Before a Golgi stack dispatches its products by budding vesicles fromt he trans face, it sorts these products and targets them for various parts of the cell. Molecular identification tags, such as phosphate groups added to the Golgi products, aid in sorting. Transport vesicles budded fromt he Golgi may have external molecules on their membranes that recognize "docking site" on the surface of specific organelles or on the plasma membrane, therefore, targeting the vesicle appropriately. Lysosomes. Lysosome is a membranous sac of hydrolytic enzymes that an animal cell uses to digest macromolecules. Lysosomal enzymes work best in the acidic environment found in lysosomes. If a lysosome breaks open or leaks its contents, the released enzymes are not very active because the cytosol has a neutral pH. However, excessive leakage from a large number of lysosomes can destroy a cell by autodigestion. Hydrolytic enzymes and lysosomal membrane are made by rough ER and then transferred to the Golgi apparatus for further processing. Proteins of the inner surface of the lysosomal membrane and the digestive enzymes are spared from destruction by having three dimensional shapes that protect vulnerable bonds from enzymatic attack. Phagocytosis is a process that amoebas and many other protists eat by engulfing smaller organisms or other food particles. The food vacuole formed , and then fuses with a lysosome and digests the food. Digestion products pass into cytosol and become nutrients for the cell. In human body, white blood cell helps defend the body by engulfing and destroying bacteria and other invaders. Lysosome use their hydrolytic enzymes to recycle the cell's own organic material; this is called autophagy. During autophagy, damaged organelle or small amount of cytosol become surrounded by a double membranes, and lysosome fuses with the outer membrane of their vesicle. The lysosomal enzymes dismantle the enclosed material, and the organic monomers are returned to the cyotosol for reuse. The lysosomes become engorged with indigestible substrates, which begin to interfere with other cellular activities. Vacuoles. Vacuoles are membrane-bounded vesicles whose functions vary in different kinds of cells. Food vacuoles are formed by phagocytosis. Many freshwater protists have contractile vacuoles that pump excess water out of the cell, thereby maintaining a suitable concentration of ions and molecules inside the cell. In plants and fungi, which lacks lysosomes, vacuoles carry out hydrolysis; The central vacuole develops by the coalescence of smaller vacuoles, themselves derived from the endoplasmic reticulum and Golgi apparatus. The vacuolar membrane is selective in transportin solutes. As result, the solution inside the central vacuole is called cell sap, is different in composition from the cytosol. It can hold reserves of important organic compounds such as proteins stockpiled in the vacuoles of storage cells in seeds. Also it is the plant cell's main repository of inorganic ions, such as potassium and chloride. Many plant cells use their vacuoles contain pigments that color the cells. Vacuoles may also help protect the plant against predators by containing compounds that are poisonous or unpalatable to animals. The vacuole has a major role in the growth of plant cells, which enlarge as their vacuoles absorb water, enabling the cell to become larger with a minimal investment in new cytoplasm. Mitochondria and Chloroplasts. Mitochondria and chloroplasts are the organelles that convert energy to forms that cells can use for work. Mitochondria are the site of cellular respiration, the metabolic process that generates ATP by extracting energy from sugars, fats, and other fuels with the help of oxygen. Chloroplasts, are found in plants and algae, and they are the sites of photosynthesis. They convert solar energy to chemical energy by absorbing sunlight and using it to drive the synthesis of organic compounds such as sugar from carbon dioxide and water. Both of them are not part of endomembrane system. Mitochondria have two membrane separating their innermost space from the cytosol, and chloroplasts have three. The membrane proteins of mitochondria and chloroplasts are made not by ribosomes bound to the ER, but by free ribosomes in the cyotosol and by ribosomes contained within these organelles themselves. They also contain small amount of DNA that programs the synthesis of the proteins made on the organelle's ribosomes. Mitochondria and chloroplasts are semiautonomous organelles that grow and reproduce within the cell. Mitochondria Mitochondria are found in all eukaryotic cells; Some cells have a single large mitochondrion, but more often a cell has hundreds or thousands of mitochondria. The number correlates with the cell's level of metabolic activity. The mitochondrion is enclosed by two membranes, each a phospholipid bilayer with a unique collection of embedded proteins. The outer membrane is smooth, but the inner membrane is convoluted, with infolding called cristae. The inner membrane divides the mitochondrion into two internal compartments : the first is the inter-membrane space, the narrow region between the inner and outer membranes and the second compartment, the mitochondrial matrix, is enclosed by the inner membrane. the matrix contains many different enzymes as well as the mitochondrial DNA and ribosomes. Enzymes in the matrix catalyze some steps of cellular respiration. Other proteins that function in respiration, including the enzyme that makes ATP are built into the inner membrane, As highly folded surface, the cristae give the inner mitochondrial membrane a large surface ares, thus enhancing the productivity of cellular respiration. Chloroplasts The chloroplast is a specialized member of related plant organelles called plastids. Chloroplasts contain the green pigment chlorophyll, along with enzymes and other molecules that function in the photosynthetic production of sugar. Its shape is lens-shaped and found in leaves. The contents of a chloroplast are partitioned from the cytosol by an envelope consisting of two membranes separated by a very narrow intermembrane space. Inside the chloroplast is another membranous system in the form of flattened, inter-connected sacs called thylakoids. Thylakoids are stacked like poker ships, and each stack is called granum. The fluid outside the thylakoids is the stroma which contains the chloroplast DNA and ribosomes as well as many enzymes. The membranes of the chloroplast divide the chloroplast space into three compartments: the intermembrane space, the stroma, and the thylakoid space. Cytoskeleton. Cytoskeleton is a network of fibers extending throughout the cytoplasm. It plays a major role in organizing the structure and activities of the cell. It is composed of three types of molecular structure: microtubules, microfilaments, and intermediate filaments. The main function of the cytoskeleton is to give mechanical support to the cell and maintain its shape. Cytoskeleton is stabilized by balance between opposing forces exerted by its elements. The cytoskeleton is more dynamic than an animal skeleton; it can be quickly dismantled in one part of the cell and reassembled in a new location, changing the shape of the cell. Also several types of cell motility involve the cytoskeleton. the cell motility encompasses both changes in cell location and more limited movements of parts of the cell. Cell motility require the interaction of the cytoskeleton with motor proteins. Cytoskeletal elements and motor proteins work together with plasma membrane molecules to allow whole cells to move along fibers outside the cell. The cytoskeleton is also involved in regulating biochemical activities in the cell in response to mechanical stimulation forces exerted by extracellular molecules via cell-surface proteins are apparently transmitted into the cell by cytoskeletal elements, and the forces may even reach the nucleus. Microtubules- thickest All eukaryotic cells have microtubules; the wall of the hollow tube is constructed from a globular protein called tubulin. Each tubulin protein is a dimer, a molecule made up of two subunits. A tubulin dimer consists of two slightly different polypeptides, alpha-tublin, and beta-tubulin. Microtubules grow in length by adding tubulin dimers. Due to the architecture of a microtubules, its two ends are slightly different; one end can accumulate or release tubulin dimers at a much higher rate than the other, therefore, growing and shrinking significantly during cellular activities. This is called the "plus end", not because it can only add tubulin proteins but because it's the end where both "on" and "off" rates are much higher. Microtubules shape and support the cell and also serve as tracks along which organelles equipped with other proteins can move. In animal cells, microtubules grow out from a centrosomes, a region that is often located near the nucleus and considered a "microtubule-organizing center". Theses microtubules function as compression-resisting girders of the cytoskeleton. Within the centrosome are a pair of centrioles, each composed of nine sets of triplet microtubules arrange in a ring. Before division, the centrioles replicate; although centrosomes with centrioles may help organize microtubule assembly in animal cell,s they are not essential for this function in all eukaryotes. Also specialized arrangement of microtubules is responsible for the beating of flagella and cilia. There are microtubule containing extensions that project from some cells. When cilia or flagella extend from cells that are held in place as part of a tissue layer, they can move fluid over the surface of the tissue. Flagella and cilia different in their beating patterns. A flagellum has an undulating motion that generates force in the same direction as the flagellum's axis. However cilia work more like oars, with alternating power and recovery strokes generating force in a direction perpendicular to the cilium's axis. A cillium may also act as a signal-receiving "antenna" for the cell. Cilia that have this function are nonmotile, and there is only one per cell. Membrane proteins on this kind of cilium transmit molecular signals from the cell's movement to its interior, triggering signaling pathways that may lead to changes int he cell's activities. Cillia-based signaling appears to be crucial to brain function and to embryonic development. Motile cilia and flagella share a common ultrastructure; each has a core of microtubules sheathed in an extension of the plasma membrane. Nine doublets of microtubules, the members of each parit sharing part of their walls, are arranged in a ring. This arrangement, referred to as the "9+2" pattern, is found in all eukaryotic flagella and motile cilia. Non-motile primary cilia have "9+0" pattern, lacking the central pari of microtubules. The microtubule assembly of a cilium or flagellum is anchored in the cell by a basal body, which is structurally very similar to a centriole. In flagella and motile cilia, flexible cross-linking proteins, evenly spaced along the length of the cilium or falgellum, connect the outer doublets to each other and to the two central microtubules Each outer doublet also has paris of protruding proteins spaced along its length and reaching toward the neighboring doublet; These are large motor proteins called dyneins, composed of several polypeptides. Dyneins are responsible for the bending movements of the organelle. A dynein molecule performs a complex cycle of movements cause by changes in these shape of the proteins, with ATP providing the energy for these changes. The mechanics of dynein-based bending involve a process that resembles walking. A typical dynein protein has two "feet" that "walk" along the microtubule of the adjacent doublet, one foot maintaining contact while the other releases and reattaches one step further along the microtubules. Without any restraints ont he movement of the microtubules doublets, one doublet would continue to "walk" along and slide past the surface of the other, elongating the cilium or flagellum rather than bending it. Microfillaments- thinnest Microfilaments are solid rods about 7nm in diameter. They are also called as actin filaments because they are build from molecules of actin, a globular protein. A microfilament si a twisted double chain of actin subunits. Microfilaments can form structural networks, due to the presence of proteins that bind along the side of an actin filament and allow a new filament to extend as a branch. The structure role of microfilaments in the cytoskeleton is to bear tension. A cortical microfilaments, a three-dimensional network formed by microfilaments just inside the plasma membrane, helps support the cell's shape. This network give the outer cytoplasmic layer of a cell called the cortex. In animal cells specialized for transporting materials across the plasma membrane, such as intestinal cells, bundles of microfilaments make up the core of microvilli. Microfillaments are well known for their role in cell motility, particularly as part of the contractile apparatus of muscle cells (myosin). Localized contraction brought about by actin and myosin also plays a role in amoeboid movement, which a cell such as an amoeba crawls along a surface by extending and flowing into cellular extension called pseudopodia. pseudopodia extend and contract through the reversible assembly of actin subunits nto microfilaments and of microfilaments into networks that convert cytoplasm fro a sol to a gel. The pseudopodium extends until the actin reassembles into a network. In plant cells, both actin-myosin interactions and sol-gel transformations brought about by actin may be involved in cytoplasmic streaming, a circular flow of cytoplasm within cells. intermediate filaments- middle range intermediate filaments are larger than the diameter of microfilaments but smaller than that of microtubules. Specialized for bearing tension (like microfilaments), intermediate filaments are a diverse class of cytoskeletal elements. Each type is constructed from a different molecular subunit such as keratins. Intermediate filaments are more permanent fixture of cells than are microfilaments and microtubules. Even after the death of the cell, intermediate filament networks often persist. Intermediate filaments are important in reinforcing the shape of a cell and fixing the position of certain organelles. For instance, the nucleus commonly sits within a cage made of intermediate filaments, fixed in location by braches of the filaments that extend into the cytoplasm. Other intermediate filaments make p the nuclear lamina that lines the interior of the nuclear envelope. In case where the shape of the entire cell is correlated with function, intermediate filaments support that shape. Cell Wall. Cell wall is an extracellular structure of plant cell that distinguishes them from animal cells. The wall protects the plant cell, maintains its shape, and prevents excessive uptake of water. The strong walls of specialized cells hold the plant up against the force of gravity. Plant cell walls are much thicker than the plasma membrane, and the exact chemical composition of the wall varies from species to species and even from one cell type to another in the same plant, but basic design of the wall is consistent. Microfibrils made of the polysaccharide cellulose are synthesized by an enzyme called cellulose synthase and secreted to the extracellular space, where they become embedded in a matrix of other polysaccharides and proteins. This combination of materials, strong fibers in a "ground substance" (matrix), is the same basic architectural design found in steel-reinforced concrete and in fiberglass. A young plant cell first secrets a thins and flexible wall called the primary cell wall; as the cell grows, the cellulose fibrils are oriented at right angels to the direction of cell expansion, possibly affecting the growth pattern. Between primary walls of adjacent cell is the middle lamella, which is a thin layer rich in sticky polysaccharides pectins. The middle lamella flues adjacent cells together. When the cell mature and stops growing, it strengthens its wall. Some plant cells do this simply by secreting hardening substances into the primary wall, but other cells add a secondary cell wall between the plasma membrane and the primary wall. Then secondary wall, often deposited in several laminated layer, has a strong and durable matrix that afford the cell protection and support. References. Berg, Jeremy M., John L. Tymoczko, and Lubert Stryer. Biochemistry. 7th ed. New York: W.H. Freeman, 2012. Print. Reece, Campbell, Lisa A. Urry, Michael L. Cain, Steven A. Wasserman, Peter V. Minosky, and Robert B. Jackson. Biology. 8th ed. San Francisco: Cummings, 2010. Print. See also. Prokaryotes 

Phospholipids | Cholesterol » Phospholipids are amphipathic molecules that are made up of a hydrophilic head and a hydrophobic tail. The head group has a 'special' region that changes between various phospholipids. This head group will differ between cell membranes [types of cells] or different concentrations of specific 'head groups'. The fatty acid tails are also differ, but there is always one saturated and one unsaturated 'leg' of the tail. Phospholipids are 2 fatty acids one saturated and one unsaturated (shown by the double bond) that are linked to a glycerol. The Noncovalent Assemblies. For the phospholipid bilayer, even though it consists of hydrophilic heads on the outer membrane, the non-covalent hydrophobic tails of the inner membrane is the key to hold the entire membrane together because there are Van der Waals attractive forces within the cell membrane in which the hydrocarbon tails are closely packed together. With its non-covalent character inside the cell membrane, the hydrophobic molecules can easily pass through the cell membrane through passive diffusion. With this, the cell has control the molecules' movement through the transmembrane proteins complexes such as pores and gates. As for the hydrophilic molecules (such as ions, carbohydrates, proteins, amino acids, and nucleic acids), they require active diffusion in order to pass through the cell membrane because of their polarity and because they are hydrophilic (most non-covalent assemblies of the cell membrane are hydrophobic, such as hydrocarbon chains). The non-covalent assemblies of the cell membrane can help give rise to bubbles such as liposome, or lipid vesicle, that can deliver drugs into a specific part of the body. The structure of liposome is very similar to the cell membrane's lipid bilayer, and the materials that compose the liposome is identical to cell membrane. Because liposome is a bubble, the structure is shaped like a ring, where the hydrophilic heads are the outer and inner ring while the hydrophobic tails are in between the hydrophilic heads' rings. Because of the ring structure, the liposomes are able to trap aqueous materials such as drugs into their rings. Once the materials are within the ring, the liposomes would deliver them to a specific location in the body, such as cancer cells. 

« Phospholipids | Cholesterol | Semi-permeability and Osmosis » Cholesterol is a major component of cell membranes and serves many other functions as well. Cholesterol helps to 'pack' phospholipids in the membranes, thus giving more rigidity to the membranes. In colder conditions cholesterol also serves to keep the fluidity in the cell membrane, by keeping space in between the phospholipids. Also cholesterol serves diverse functions such as: it is converted to vitamin D (if irradiated with Ultra Violet light, modified to form steroid hormones, and is modified to bile acids to digest fats. 

« Cholesterol | Semi-permeability and Osmosis | Proteins and channels » The membranes of cells are a fluid, they are semi-permeable, which means some things can pass through the membrane through osmosis or diffusion. The rate of diffusion will vary depending on the its: size, polarity, charge and concentration on the inside of the membrane versus the concentration on the outside of the membrane. When something is permeable it means that something can spread throughout, like (The perfume is permeating the room.). Here is a list of some molecules and how they relate to passing through the membrane without assistance, in other words, through diffusion: Ions. Various substances will pass through the membranes at varying rates through diffusion. 

Introduction. The number of known organic compounds is quite large. In fact, there are many times more organic compounds known than all the other (inorganic) compounds discovered so far, about 7 million organic compounds in total. Fortunately, organic chemicals consist of a relatively few similar parts, combined in different ways, that allow us to predict how a compound we have never seen before may react, by comparing how other molecules containing the same types of parts are known to react. These parts of organic molecules are called functional groups. The identification of functional groups and the ability to predict reactivity based on functional group properties is one of the cornerstones of organic chemistry. Functional groups are specific atoms, ions, or groups of atoms having consistent properties. A functional group makes up part of a larger molecule. For example, -OH, the hydroxyl group that characterizes alcohols, is an oxygen with a hydrogen attached. It could be found on any number of different molecules. Just as elements have distinctive properties, functional groups have characteristic chemistries. An -OH group on one molecule will tend to react similarly, although perhaps not identically, to an -OH on another molecule. Organic reactions usually take place at the functional group, so learning about the reactivities of functional groups will prepare you to understand many other things about organic chemistry. Memorizing Functional Groups. Don't assume that you can simply skim over the functional groups and move on. As you proceed through the text, the writing will be in terms of functional groups. It will be assumed that the student is familiar with most of the ones in the tables below. It's simply impossible to discuss chemistry without knowing the "lingo". It's like trying to learn French without first learning the meaning of some of the words. One of the easiest ways to learn functional groups is by making flash cards. Get a pack of index cards and write the name of the functional group on one side, and draw its chemical representation on the other. For now, a list of the most important ones you should know is provided here. Your initial set of cards should include, at the very least: Alkene, Alkyne, Alkyl halide (or Haloalkane), Alcohol, Aldehyde, Ketone, Carboxylic Acid, Acyl Chloride (or Acid Chloride), Ester, Ether, Amine, Sulfide, and Thiol. After you've learned all these, add a couple more cards and learn those. Then add a few more and learn those. Every functional group below is eventually discussed at one point or another in the book. But the above list will give you what you need to continue on. And don't just look at the cards. Say and write the names and draw the structures. To test yourself, try going through your cards and looking at the names and then drawing their structure on a sheet of paper. Then try going through and looking at the structures and naming them. Writing is a good technique to help you memorize, because it is more active than simply reading. Once you have the minimal list above memorized backwards and forwards, you're ready to move on. But don't stop learning the groups. If you choose to move on without learning the "lingo", then you're not going to understand the language of the chapters to come. Again, using the French analogy, it's like trying to ignore learning the vocabulary and then picking up a novel in French and expecting to be able to read it. Functional groups containing .... In organic chemistry functional groups are submolecular structural motifs, characterized by specific elemental composition and connectivity, that confer reactivity upon the molecule that contains them. Common functional groups include: "Note: The table above is adapted from the Functional Groups table on Wikipedia." Combining the names of functional groups with the names of the parent alkanes generates a powerful systematic nomenclature for naming organic compounds. The non-hydrogen atoms of functional groups are always associated with each by covalent bonds, as well as with the rest of the molecule. When the group of atoms is associated with the rest of the molecule primarily by ionic forces, the group is referred to more properly as a polyatomic ion or complex ion. And all of these are called radicals, by a meaning of the term "radical" that predates the free radical. The first carbon after the carbon that attaches to the functional group is called the alpha carbon. Alcohol containing two hydroxyl groups are called glycols. They have both common names and IUPAC names. Mnemonics for Functional Groups. These are possible mnemonics for the common functional groups. Vowels: Remember the vowels "A", "E", and "Y" for Alkane, Alkene, and Alkyne. Alkanes have only single covalent bonds. Alkenes have at least one double bond. Alkynes have at least one triple bond. The letters "I", "O", and "U" are not used. Furthermore, "O" and "U" would result in awkward pronunciations. Alcohol: Look for the "C-O-H" in "Alcohol." Ether: Ethers were anesthetics used in the 1800s. Dr. Kellogg also lived at the same time. Corn Flakes are made by Kellogg's. A rooster or cock (C-O-C) is the cornflake mascot. Or, think there is a C on "either" side of the O. Amine: Remember the "N" stands for nitrogen. Aldehyde: This sounds like "Adelaide," the Australian city. Australia is at the end of the Asian islands, and aldehydes are at the end of the hydrocarbon chain. The "Y" indicates a C=O double bond. Ketone: Imagine the diagonal strokes of "K" forming the C=O double bond. Carboxylic Acid: "Box" stands for boxed wine or C-O-H, alcohol. The "Y" indicates a C=O double bond. Ester: This sounds like "Estelle" George Costanza's mother in the TV show Seinfeld. George's nickname was Koko or Coco. So think of O=C-O-C. Amide: Amine with a "D". D for double. 

This book is on "abstract algebra" (abstract algebraic systems), an advanced set of topics related to algebra, including groups, rings, ideals, fields, and more. Readers of this book are expected to have read and understood the information presented in the Linear Algebra book, or an equivalent alternative. Information for contributors. This wikibook shall give an introduction to the fundamental concepts of abstract algebra, such as groups, rings and ideals, and fields and Galois theory. Contents. Fields. /Sources/ 

Hesitation noises. Hesitation noises, or vocal pauses, are the "uh" and "um" of a language, filler sounds we produce when we pause to think while speaking. Using the correct hesitation noises can make the difference between natural sounding Japanese, and awkward foreign sounding speech. They are: 

Japanese is spoken by 130 million people. This makes it the ninth most spoken language by native speakers. Linguists debate over the classification of the Japanese language, and one general theory asserts that Japanese is an isolated language and thus a language family of its own, known as Japonic languages. Another major theory includes Japanese as part of a hypothetical Altaic language family which spans most of Central Asia and would also include Turkic, Mongolic, Tungusic, and Korean languages. Neither of these theories has yet been generally accepted. Japan is the only country where Japanese is the sole official language (though the island of Angaur has Japanese as one of three official languages). There are, however, numerous speakers in other countries. These are largely due to emigration, most notably to the United States of America (California and Hawaii, in particular), Brazil and the Philippines. Furthermore, when Japan occupied and colonized much of East Asia, Southeast Asia, and the Pacific, the locals were educated in the Japanese language. Many elderly locals in Korea, Taiwan, and parts of China still speak Japanese. Japan has steadily developed for many centuries and, unlike many other cultures, has not been seriously affected by any major invasions until recent times. A substantial part of the vocabulary, though, has been borrowed over the years from Chinese, Portuguese, Dutch, German, French, and most recently English. Grammar. While Japanese grammar is very regular, it is markedly different from English. Japanese has been deemed a subject–object–verb (SOV) and topic-prominent language, whereas English is a subject–verb–object (SVO) and subject-prominent language. To illustrate, the English sentence “Cats eat mice” contains a subject (cats), a verb (eat), and an object (mice), in an SVO order, where the “-s” is a "plural marker", and “mouse” → “mice” is a plural marker by ablaut, but only the word order indicates which is the subject and the object—i.e. which is dining and which is the meal. The topic-prominence is not obvious in this example; “cat” is the subject (the agent) in English, but it is the topic (what the sentence is "about") in Japanese. In the above example, “は” is the topic, and “を” is the comment. The verb “kū” means “eat” in the sense of one animal consuming another. To speak about a person eating, it would make more sense to use the word “taberu” which means “eat,” as in to consume a meal. Japanese does not have articles (the words, “a” or “an”, or “the”), nor is it mandatory to indicate number (singular versus plural). In the sentence above, “” could mean either “cat” or “cats.” The “mouse/mice” ablaut does not occur in Japanese, which is an agglutinative language (inflecting by appending) and highly regular. In Japanese, the plural is formed by adding the ending “-tachi,” or “-ra.” Thus, the word, “cats” would be “nekotachi,” and is always plural, but the word “neko” can be either singular or plural. “Nekora,” however, would sound rather strange, as “-ra” and “-tachi” are not necessarily interchangeable. For the beginner, therefore, it is best not to worry about learning plural endings. For the English speaking student of Japanese grammar, the greatest hurdles to cross are probably the thought process of the Japanese sentence and learning the seemingly endless variety of endings available for modifying verbs and the order in which they can be strung together. The grammatical paradigm of SVO or SOV is completely irrelevant in the study of Japanese and other languages outside the Indo-European family of languages. In truth, it is not only unimportant, it is untrue, and will cause the student of the language to fail in acquiring fluency, because it is an artificial imposition of an Indo-European construct on a non-Indo-European language. Japanese, like Tagalog and many other languages, uses affixes to explicitly demonstrate grammatical relationships instead of using syntax. In Japanese, word order will not change the meaning of the sentence. However, it will change the emotional character. An SVO word order is not incorrect in Japanese, and native speakers use it frequently, as a matter of fact to heighten the emotional charge. Thus, “あれは何だ” is a simple question: “What is that?” But “何だあれは” should receive an exclamation point at the end because the word order indicates that the speaker is clearly upset or at least annoyed by whatever “that” is. Japanese sentences thus are not SOV. They are TV: T stands for topic and V for verb. Verbs are really the secret to success in acquiring fluency in Japanese. Thus, greatest attention should be given to learning the verb forms. There are two tenses of time: past and present. The present tense is used to describe future events. All past tense verbs have the ending “-た” (“-ta”) or “-だ” (“-da.”) The present tense always ends in the vowel “-u” in the positive and “-nai” in the negative. There is only one exception: the word, “だ” (“da”), which is the present tense of “be” (“am,” “are,” “is.”) As in probably all languages, this verb is highly irregular in Japanese and its usage must simply be memorized. For the English speaker the two time tenses should be quite easy to remember because in English the past tense is usually indicated by a final “-t” or “-d,” and the present tense of the basic positive present tense verb “do” ends in the “u” sound. The “-nai” ending sounds similar to “nay” in English. In conversational Japanese, a complete sentence will end in a present tense or past tense verb. As indicated earlier, there are many other possible endings, but they are not used in the final position to complete a sentence, nor are they used at the end of the active verb. Beyond the verb, there are words that indicate the function of words and phrases as they relate to the verb. The most important of these are “は” (“wa”), “が” (“ga”), “に” (“ni”), “の” (“no”), “を” (“o”), and “で” (“de.”) “は” marks what is being discussed. “が” follows the word that is the agent of the verb. This means who is doing something is the active tense, and who is receiving the action of a passive verb. In both cases, they mark “the who.” “に” indicates the direction toward and is usually translated as “in,” “to,” “at.” “の” indicates possession or source and usually is translated as “-’s” or “of.” “を” is only appropriate with active transitive verbs because it marks the direct object. Finally, “で” at the end of a place word indicates where a verb happened. It is usually translated “at,” “on,” or “in.” Added to the end of word that represents an object, it marks the instrument of the verb, what was used to perform the verb. It is translated, “with.” The example translates to "Speaking of cats, Pitchan came home and ate the dog’s food in the kitchen." (The “て” verb ending indicates incompletion.) Thus, every word or phrase in a Japanese sentence takes an ending that explicitly denotes the function of that word or phrase and how it relates to the verb. Levels of politeness. Japanese culture and society is based on a hierarchy of higher status (目上 "meue") and lower status (目下 "meshita"). As such, there are three varying levels of politeness. Because Japanese is primarily a "vertical" society, all relationships contain an element of relative station. For example, a student is a lower station than a teacher, and therefore a student would use polite language when speaking to a teacher, but the teacher would use plain language when speaking to a student. A salesperson talking to a customer would place himself/herself far below the customer, and would therefore use honorific language, whereas the customer would use either plain or polite language. Honorific language is not a separate category from plain and polite language, but a separate concept that uses different rules. When using honorific language, a Japanese speaker modifies nouns, verbs, and adjectives to either lower himself/herself and their associates, or exalt someone else and that individual's associates. Whereas the use of plain or polite language is determined by the relative station of the person "to" whom you are speaking, the use of honorific language is determined by the relative station of the person "about" whom you are speaking. Exalted language is applied when you are speaking about someone who is due respect, such as a professor, an executive, a political official, or a customer. Exalted language is only applied to other people, never to oneself. Humble language, however, is "only" applied to oneself and people associated with oneself. It would be inappropriate, for example, to use humble language to describe a beggar, even though they would be extremely low on the social ladder. The Japanese writing system. Japanese is written mostly using three writing scripts, "kanji", "hiragana" and "katakana". Kanji are Chinese characters that were first introduced to Japan in the 4th century. Unlike Chinese, Japanese is a highly inflected language with words changing their ending depending on case, number, etc. For this reason, the hiragana and katakana syllabaries were created. The hiragana serve largely to show the inflection of words, as conjunctions and such. The katakana are mainly used for loan-words from other languages. Kanji. The Japanese writing system is derived from the Chinese ideographic character set (Japanese: 漢字 "kanji", Mandarin: 汉字 "hanzi"). They are usually very similar to Traditional Chinese characters. Though kanji are Chinese in origin their use is dictated by Japanese grammar. Each character may be read in different ways depending on the context it is in. The number of existing Chinese characters has been variously estimated at between 40,000 and 80,000; however, only a small subset is commonly used in modern Japanese. An educated Japanese person will generally be able to read between 2,000 and 4,000 characters. In order to be literate in the Japanese language, the student should strive to master at least the 2,136 general-use characters (常用漢字 – "jōyō kanji") established by the Ministry of Education. Hiragana and katakana. The syllabaries, known as "kana" (), were developed around 900 AD by simplifying kanji to form the "hiragana" (ひらがな, or 平仮名) and the more angular "katakana" (カタカナ, or 片仮名). Hiragana can be recognized by the characteristic curved shapes, while katakana are identifiable by their sharp edges and straight lines. The creation of one of the scripts has been attributed to Kūkai (774-835, alias Kōbō Daishi) the famous monk who introduced Shingon Buddhism to Japan. Hiragana and katakana are almost completely phonetic—much more so than the English alphabet. Each set, however, is referred to as a "syllabary" rather than an alphabet because each character represents an entire syllable with only a single consonant (which is a more recent addition) (see ../Pronunciation/ for more). The syllabary charts in Japanese are referred to as the "gojūon" (), meaning "fifty sounds" because they are written in a five by ten chart. However, there are a few gaps in the table where certain sounds have fallen out of use. Modern Japanese can be written using 46 kana. In practical use, hiragana is used to write, for example, inflectional endings for adjectives and verbs (送り仮名 "okurigana"), grammatical particles (助詞 "joshi") and auxiliaries (助動詞 "jodōshi"), Japanese words that have no kanji (or not commonly known kanji), and annotations to kanji to indicate pronunciation (振り仮名 "furigana"). Katakana is used to write, for example, foreign words and names, onomatopoeia, emphasized words (somewhat like italicized words in English text), and technical and scientific words, such as plant, animal, and mineral names. Contents 

Uniform Resource Identifier (URI) The URI is a phase that identifies any object on the internet by specifying the protocol, web address, file location, and name of the object. The URI may also contain data for delivery to the target server, such as a user id or query. A "URL" is one form of a URI. When accessing the Web, The URI is usually displayed in a special field in the "toolbar" of a "browser". URI originally stood for a "Universal Resource Identifier", and details are available from the "Internet Engineering Task Force" in RFC1630 at IETF Web site. The following regex may be used to validate strings for the RFC2396 specification: /^(https?|ftp):\/\/(?# protocol )(([a-z0-9$_\.\+!\*\'\(\),;\?&amp;=-]|%[0-9a-f]{2})+(?# username )(:([a-z0-9$_\.\+!\*\'\(\),;\?&amp;=-]|%[0-9a-f]{2})+)?(?# password )@)?(?# auth requires @ )((([a-z0-9][a-z0-9-]*[a-z0-9]\.)*(?# domain segments AND )[a-z]{2}[a-z0-9-]*[a-z0-9](?# top level domain OR )|(\d|[1-9]\d|1\d{2}|2[0-4][0-9]|25[0-5]\.){3}(?#  )(\d|[1-9]\d|1\d{2}|2[0-4][0-9]|25[0-5])(?# IP address ))(:\d+)?(?# port ))(((\/+([a-z0-9$_\.\+!\*\'\(\),;:@&amp;=-]|%[0-9a-f]{2})*)*(?# path )(\?([a-z0-9$_\.\+!\*\'\(\),;:@&amp;=-]|%[0-9a-f]{2})*)(?# query string )?)?)?(?# path and query string optional )(#([a-z0-9$_\.\+!\*\'\(\),;:@&amp;=-]|%[0-9a-f]{2})*)?(?# fragment )$/i 

Japanese verbs, (動詞; どうし), inflect heavily to indicate "formality", "tense" or "mood", primarily in their ending. There are two tenses, several levels of formality and three classes of verbs, depending on their inflection. The two tenses are perfective (often considered past tense) and present (or technically, non-past, as the future tense is not indicated). Out of the several levels of formality, two are the most common: plain and polite. Japanese verbs are officially categorised into five classes, but as two of these inflect much the same and another two only contain one verb each, these are usually merged into three when Japanese is taught as a foreign language. These are the "consonant stem"-, "vowel stem"- and "irregular" classes. Dictionaries use the plain present positive form (commonly known as "dictionary form") as the headword for verbs. Verbs are classed based on their conjugations. Their endings don't determine the class, but are a general indicator. Different inflections can also have suffixes. These may also be verbs with their own conjugations. Not all suffixes can be used on all verb inflections and others may only follow the verb stem. Examples are conjunctive + いる, せる・させる (causative), and られる (potential). Ignoring the formality and the negative conjugations, the following is a list of verb conjugations Ichidan class. Vowel-stem verbs end on a full syllable (hence the term: "vowel"-stem). In a sense, the final "る" of the dictionary form is dropped and the respective endings just added on. The Japanese term "" refers to the fact that the stem ending occupies only one row in the kana chart. The following table shows a few forms of the verb "食べる" (たべる, "e." to eat): Godan class. Consonant-stem verbs end in the middle of a syllable (hence the term; "consonant"-verb). That syllable changes depending on the form. The plain form has an "u" sound ("u", "tsu", "ru", "ku", "gu", "bu", "mu", "su"), the "-masu" form has an "i" sound ("i", "chi", "ri", "ki", "gi", "bi", "mi", "shi"), and the negative form has an "a" sound ("wa", "ta", "ra", "ka", "ga", "ba", "ma", "sa"). The potential form has an "e" sound ("e", "te", "re", "ke", "ge", "be", "me", "se") and the volitional form has an "ō" sound ("ō", "tō", "rō", "kō", "gō", "bō", "mō", "sō"), so putting these together with the sounds above shows that verb conjugations follow the vowel syllabary of the Japanese character set:　あ "a", い "i", う "u", え "e" and お "o". The Japanese term "" comes from the fact that the stem's last syllable spans all five rows of the kana chart in at least one form. The following table shows a few forms of the consonant-stem verb "話す" (はなす "e." to speak). The て-form (conjunctive) and past positive form of a consonant-stem verb change the root for euphony according to the last syllable of the root (example in parentheses): 行く (いく) (to go) has an exceptional て-form 行って (いって). If the verb stem ends on "う" such as in the verb 買う(かう, "e." to buy) then its negative stem becomes -わ as in 買わない ("to not buy"). This is because the root is treated as kawu (despite the "wu" syllable not existing in modern Japanese). Irregular verbs. Two common verbs do not share a conjugation pattern with any other verb. They are therefore commonly classed as "irregular" verbs. Formally, they are called "変格" (へんかく) verbs, as opposed to the regular "正格" (せいかく) verbs. This construction is made to use verbs and nouns of Chinese origin, for instance, from Chinese "確認"　("què rèn", confirmation) is formed in Japanese the verb "確認する" (かくにんする), or "約分" ("yuē fēn", simplify a fraction (math.)) which derives into "約分する"　(やくぶんする). The forms are "する" ("e." to do, as in the examples) and "" ("e." to come). The following table shows some of their conjugation forms. Many verbs end on "〜する" and can be grouped in three categories: The only commonly-used combination with "来る" is "やってくる", meaning "to come". Polite forms. The polite (or formal) forms are simple as all of the consonant-stem verbs sit in the い-line (行く→行き) and the inflections are the same for consonant- and vowel-stem verbs. The following table shows the polite forms for "行く" (いく, "e." to go): The imperative (〜ませ) is not used in formal forms except for a few polite verbs (see below). Other irregularities. A small number of verbs tend to be conjugated differently from the groups that they are normally placed in. Polite language. The verbs below are all consonant stem verbs but conjugate differently. While the regular forms also exist, they are seldom used. The conjunctive and past forms of the first two verbs, "くださる" and "なさる", also have the alternative forms "くだすって／くだすった" and "なすって／なすった", in addition to the normal regular conjugations "くださって／くださった" and "なさって／なさった". These alternative forms have, however, fallen into disuse. While they are often encountered when reading texts from a few decades ago, the regular conjugations are usually used today. The first three of the above verbs are also the only ones where the imperative form "ませ" of the auxiliary verb, "ます", is used to add an extra level of politeness: Additionally, ございます, which originally came from the now-defunct "yodan" (四段, "e." four-row) classical Japanese verb "ござる", is also used, although in modern usage, it is always used with the ます auxiliary verb ending. There is no imperative form (i.e. you cannot use ませ like above). 得る. 得る (うる/える, "e." to get, or to be able to) is the only surviving "nidan" (二段, "e." two-row) class verb in modern Japanese. It has conjugations as in the below table: "得る" can be read both as "える" in its "terminal form" (at the end of the sentence, or in situations such as attaching to べき). The "うる" reading is also used in those situations and in the "attributive form" (e.g. when attached to nouns). It is therefore incorrect to say "えるもの" as the correct form would be "うるもの". The combination "あり得る" is normally read "ありうる" in the present forms. All other conjugations follow the table above. Miscellaneous irregularities. The vowel stem verb "呉れる" (くれる "e." ) imperative form "くれ" (rather than the expected "くれろ"). Other "くれる" verbs of other unrelated meanings conjugate to the usual "くれろ". The consonant stem verb "ある" expresses existence, but absence is expressed with the adjective "ない". Note that many textbooks also treat "ない" as a verb. The reader may also wish to be aware that more formal "ぬ" negative form and its conjunctive form, "ず", are still used: "あらぬ"/"あらず". Summary of verb conjugations. See the Wikipedia page for present negative, past and past negative forms of "i" and "na" adjectives. 

Internet Society (ISOC) The ISOC "(pronounced eye-sok)" is a cooperative of commercial, professional and governmental organizations that promotes internet technologies. It was founded in 1991 as a corporate umbrella for the "Internet Engineering Task Force (IETF)" and other ad-hoc activities ISOC now sponsors an annual conference and other events. It is headquartered in Virginia and Geneva, and its web address is www.isoc.org. 

^ Esperanto ^ | About the book: Authors The Esperanto textbook was started by . Other contributors have included: ^ Esperanto ^ | About the book: Authors 



G: The French alphabet. In addition, French uses several accents which are worth understanding. These are: à, è, ù, (grave accents) and é (acute accent). A circumflex applies to all vowels: â, ê, î, ô, û. A tréma (French for dieresis) is also applied: ä, ë, ï, ö, ü, ÿ. Two combined letters are used: æ and œ, and a cedilla is used on the c to make it sound like an English s: ç. V: Time. In French, “il est” is used to express the time; though it would literally translate as “he is”, it is actually, in this case, equivalent to “it is” (unpersonal "il"). Unlike in English, it is always important to use “heures” (“hours”) when referring to the time. In English, it is OK to say, “It’s nine,” but this wouldn’t make sense in French. The French time system traditionally uses a 24-hour scale. Shorthand for writing times in French follows the format "17h30", which would represent 5:30PM in English. V: The days of the week.. Notes: V: Seasons.  "Bien..." is an adverb meaning "well". Its adjective equivalent is "bon(ne)", which means "good". Since "je vais", meaning "I go", uses an action verb,&lt;br&gt; the adverb "bien" is used. In English, I'm good, which uses the linking verb "am", is followed by an adjective rather than an adverb.  "Est-ce que..." doesn't mean anything (like the Spanish upside down question mark) and is used to start a question.&lt;br&gt; This can be used in a similar manner to "do" in English. Instead of "You want it?", one can say "Do you want it?"  "chez..." is a preposition meaning "at the house of...". "Chez moi" is used to say "at my place". "Chez ["name"]" is used to say "at ["name"'s] place".  "on..." can mean "we" or "one". 

Alphabet. The Esperanto alphabet has 28 letters. Four letters from the English alphabet have been dropped – Q, W, X and Y – and there are six new accented letters: Ĉ, Ĝ, Ĥ, Ĵ, Ŝ and Ŭ. The first five have an angle-shape accent called a "circumflex" (^) over them, whilst the last has an accent rather like the bottom part of a circle, which is called a "breve" (˘). All of the accented letters are unique to Esperanto except for "ŭo" (Ŭ), which also exists in Belarusian, and "ĝo" (Ĝ), which also exists in Aleut. Some of the accented letters may be used in transcription systems for languages that use non-Latin alphabets. Vowels. As in English, five letters are vowels (A, E, I, O, U), and the rest are consonants. The letter "ŭo" (Ŭ) is a consonant, not a vowel. Collation. Collation in Esperanto is the same as for English, except that the accented characters are counted as separate characters and collated after their non-accented versions. Collation is as shown in the table above. Pronunciation. Each letter in Esperanto has only one pronunciation (allowing for cultural variation), and no letters are silent. There are six dipthongs (see the next section), but their pronunciation follows logically from their constituent letters, except for being shortened into a single syllable. This means that Esperanto is pronounced just as it is spelled. Also, each sound has only one way of being written, so it is very easy to spell Esperanto words you hear. The technical description for these traits is that Esperanto is "phonetic" and "orthographic". Stress. The stress on every word is put on the penultimate (second-to-last) syllable. 

Be sure to make use of the pronunciation appendix. Grammar. Since Esperanto is a very regular language, all its rules can be applied universally without exceptions. This means Esperanto grammar concepts are much easier to understand than those of natural languages. In this first lesson, we shall examine how to form nouns, adjectives, and the present tense. Nouns. In any language, "nouns" are words that designate a person, place, thing, idea, or quality. Some examples of nouns in English are: "house", "friends", "cake", "John", "France", and "gardens". In Esperanto, all nouns end in -o. The part of the word that goes before the -o is known as the "root". For example, in the word urbo (city), urb- is the root and the -o makes it a noun. To make a noun plural, add a -j to the end, for example urboj (cities). Some examples of nouns in Esperanto: (human), (house), (friends), (cake), (John), (France) and (gardens). To say "a" or "an", as in "a town", just say the noun on its own, e.g. urbo (town, a town). There is no indefinite article ("a" or "an") in Esperanto. The word for "the" is la, e.g. la urbo (the city). La never changes for singular or plural. Adjectives. "Adjectives" are words that describe a noun. Some English examples are: "happy", "tired", "beautiful", "young" and "fresh". To change an Esperanto noun into its corresponding adjective, replace the -o with an -a. For example, urbo (town) gives rise to urba (urban, "relating to a town"). Some examples of adjectives in Esperanto: (happy), (tired), (beautiful), (young) and (fresh). In Esperanto, an adjective must "agree in number" with the noun it describes. This means that if the noun is singular, the adjective must also be. If the noun is plural, the adjective must also be, too. Some examples: la freŝa kuko (the fresh cake), la freŝaj kukoj (the fresh cakes); feliĉa homo (a happy person), feliĉaj homoj (happy people). The prefix mal- changes an Esperanto word into its opposite meaning – a feature that greatly reduces the vocabulary. Here are some examples of mal- words in Esperanto: (unhappy), (alert, not tired), (ugly), (old) and (stale). Adverbs. "Adverbs" are words that describe a verb, an adjective, or another adverb. They indicate manner, place, time or quantity. Some English examples are: "quickly", "orally", "at home", and "in writing". To change an Esperanto word into an adverb, replace the usual ending (-a for adjectives, -o for nouns, and -i for verbs) with -e. The meaning of the base word determines whether it becomes a manner, place, time or quantity adverb. Some examples of adverbs in Esperanto: (quickly), (orally), (at home), and (in writing). Please note that not all adverbs use this rule, but the overwhelming majority of them do. Once you are introduced to these adverbs, it will be obvious why there are exceptions. Personal pronouns. In Esperanto there are ten personal pronouns. However, you will initially need to only know seven of these pronouns. Possessive pronouns. To turn a personal pronoun into a possessive pronoun (which is an adjective), simply add an -a to the end. Verbs – present tense. The basic form of a verb is called its "infinitive". In English, this is the part of the verb that has "to" in front of it, as in the sentence "John likes to play football". In Esperanto, the infinitive simply has an -i after the root, e.g. (a game), (to play). The present tense has three forms in English. For example, one can say either "I kick", "I am kicking" or "I do kick"; "he laughs", "he is laughing", "he does laugh"; "Robert eats the cake", "Robert is eating the cake", "Robert does eat cake". All these forms are represented under one form in Esperanto. To form the present tense of any Esperanto verb, simply substitute the -i in the root with -as. Some examples: mi legas (I am reading), li ridas (He is laughing), Roberto manĝas la kukon (Robert eats the cake). You will also note that there is no verb conjugation in Esperanto for person. For example, in English you would say "Bob eats", "I eat", and "She eats"; In Esperanto, you would use the same form of the verb: "Bob manĝas", "Mi manĝas", "Ŝi manĝas". Objects and the accusative case. Like in English, every complete declarative sentence in Esperanto requires at least two parts: a subject and a verb. In the sentence "I ate", the subject is "I" and the verb is "ate". A subject is a noun which performs an action. However, in the sentence "I ate spaghetti", there is a third word: "spaghetti". The word "spaghetti" in this sentence is what is known as a direct object. A direct object is a noun which is having an action performed on it. It is being "verb'ed", so to speak. Take special note of the last sentence. In the first three sentences, the direct object directly followed the verb. However, in the last sentence, the word "Jane" directly follows the verb "gave". So why isn't Jane the direct object of the sentence? Because Jane is not having an action performed on her by the verb of the sentence, "give". Frank is not giving Jane, Frank is giving flowers "to Jane". Since the flowers are what is being given, the flowers are the direct object. So what is Jane in this sentence, then? Jane is what is known as an indirect object. An indirect object is a noun which is neither performing an action nor having an action performed directly upon it, but receiving the action of the verb less directly than the "direct object" (hence the name). In Esperanto, the indirect object will always take a preposition. For example, "Frank baked Jane a cake" becomes "Frank baked a cake for Jane." This basal reader stuff is real cute. You'd almost believe that English could be a clear means of communicating. However, something working in the lab on small toy problems doesn't mean it will work on the complexities of real life. So now let's try some more advanced sentences in English. As you can see, English sometimes chokes on the complexities of real life. The incorrect understanding of the sentence about the First Amendment has the same basic structure as "The pilot makes the heading in the flight plan to the north.", giving the false parse a similar structure to "The amendment redirects the trust in their ability to criticize, into the Constitution.", which would mean "The First Amendment has persuaded the citizenry to trust that they are able to criticize the governor's administration of the presumptive draft on the grounds that the governor is administrating the presumptive draft unconstitutionally instead of trusting that they are able to criticize the governor's administration of the presumptive draft on other grounds.", which, depending on the minute details of history, the English language, and then-current law, could actually make sense, but is not the correct understanding of the grammar or the meaning of the sentence. People can sort of use their sapience to get past part of this problem, but in controversies, debates, and trials where words are deliberately or accidentally twisted, in technical fields, in natural language processing, and with English as a Foreign Language, this obviously isn't working. Now you know one of the many reasons why natural language processing is an unsolved problem in computer programming. Another reason is that most English words have multiple meanings, up to 60 or even 100 meanings. Most Esperanto words only have one meaning, and its words with more than one meaning usually only have 2 or 3 meanings. Google Translate makes countless errors, and as such is no substitute for taking the courses over there at Lernu. We'll try that again by shoehorning it into Esperanto. These sentences are kind of stilted, and in another situation one would use a very different wording. The stilted wording more or less preserves the basic structure of the English sentences, though. We'll start with the translation of "Now let's try more advanced sentences in Esperanto.", which is "Nun ni provu iom Esperantajn frazojn kiuj estas pli progresintaj." In English, word order in a sentence helps determine whether a noun is a subject, a direct object, or an indirect object. In English, the order is often "subject, verb, direct object", "direct object, subject, verb" or "subject, verb, indirect object, direct object", and various other orders. As you can see from the table above, there are many plausible but incorrect comprehensions of the English sentences, but no plausible incorrect comprehensions of their Esperanto translations. That's because grammar in planned languages, such as Esperanto, tends to be much more precise than in natural languages. Esperanto uses affixes to explicitly and unambiguously mark words for part of speech, case, gender and number. Thus, part of speech tagging of Esperanto's single-stem words is already solved. In order to make a direct object in Esperanto, one simply adds the letter "n" to the noun (if the noun is plural, the "n" is added after the "j"). Therefore, in Esperanto, subjects, verbs, and direct objects can be put in any order, with little, if any, loss of clarity. All of the following sentences, which mean "the apple loves the banana" are grammatically correct in Esperanto. Conversation – Introducing yourself. Two people – Jean and Frank – meet for the first time: 



Audio Help 

^Lesson 3^ "Fill in the blanks". Por favor, rellena los espacios en blanco con la forma correcta del verbo: "Please, fill in the blanks with the correct form of the verb:" 1. Nosotros _________ (aprender) español. "We learn Spanish." 2. Yo _________ (comprar) un libro. "I buy a book." 3. Carmen y Roberto _________ (viajar) a Mexico. "Carmen and Roberto travel to Mexico." 4. Ana __________ (hablar) ingles. "Ana speaks English." 5. Tú ________ (beber) una cerveza. "You drink a beer." 6. Susana ________ (escribir) una carta. "Susana writes a letter." 7. Los niños _________ (estudiar) para el examen. "The children study for the exam." 8. Fernando y Lucas __________ (cantar) una cancion. "Fernando and Lucas sing a song." 9. Tú ________ (leer) un libro. "You read a book." 10. Vosotros _________ (subir) las escaleras. "You all (plural) climb the stairs." Soluciones a los ejercicios "Solution to exercices" ^Lesson 3^ 

This page is a list of links to the solutions to the exercises in this Wikibook. The exercises themselves can be found here. 

Ruby is an interpreted, object-oriented programming language. Its creator, , aka “Matz”, released it to the public in 1995. Its history is covered here. Its many features are listed here. The book is currently broken down into several sections and is intended to be read sequentially. Getting started will show how to install and get started with Ruby in your environment. Basic Ruby demonstrates the main features of the language syntax. The Ruby language section is organized like a reference to the language. Available modules covers some of the standard library. Intermediate Ruby covers a selection of slightly more advanced topics. Each section is designed to be self contained. Table of Contents. Ruby Semantic reference. See also some rdoc documentation on the various keywords. Built in Classes. This is a list of classes that are available to you by default in Ruby. They are pre-defined in “core.” Available Standard Library Modules. These are parts of Ruby that you have available (in the standard library, or via installation as a gem). To use them you typically have to require some filename, for example codice_62 would make accessible to you the Tracer class. You can see a list of basically all the (std lib ruby) modules available in the ruby source and lib readme. There are a several more modules available in the std lib, which are C based extensions. You can see their list here. Other Libraries. GUI Libraries. Here is info on some specifically: Intermediate Ruby. Here are some more in depth tutorials of certain aspects of Ruby. 

^Lesson 1^ Pronombres personales. "(Personal pronouns.)" Ejercicio 1. Rellena los espacios en blanco con los pronombres personales correctos en español. "Exercise 1." "Fill the blank spaces with the correct Spanish personal pronouns." Ejercicio 2. Rellena los espacios en blanco con los pronombres personales correctos en español. "Exercise 2." "Fill the blank spaces with the correct Spanish personal pronouns." Soluciones a los ejercicios "Solution to exercices" ^Lesson 1^ 

^Lesson 1^ Llena los espacios en blanco con la forma verbal correcta del verbo "ser". "Fill the blank spaces with the correct form of the verb "ser" (to be)." Llena los espacios en blanco con la forma verbal correcta del verbo "estar". "Fill the blank spaces with the correct form of the verb "estar" (to be)." Llena los espacios en blanco con la forma verbal correcta del verbo "ser" o "estar". "Fill the blank spaces with the correct form of the verb "ser" or "estar" (to be)." Soluciones a los ejercicios "Solution to exercices" ^Lesson 1^ 

Cascading Style Sheet (CSS) is an HTML feature using a set of hierarchic input files that specify display characteristics used to interpret various "HTML" formatting tags. The CSS2 version also supports the description of display properties that are applied to "XML" applications. The CSS is an open specification maintained by the "W3C", and detailed online information describing style sheets may be found at the W3C's Specification Page. 

Information on the Oxford, Cambridge and RSA examinations is available from the OCR Website. End of Module tests. See also Electronics. 

GCSE Science/Electricity Have you ever noticed a crackle when you pull off a sweater? Have you ever had your hair stick to your face after combing? What about seeing a balloon stick to the wall after it has been rubbed? All these effects are caused by static electricity. Static electricity is the imbalanced charges of matter. It was the Greeks who first noticed that when they rubbed a piece of amber, it attracted small pieces of paper. The Greek word for amber is elektron, and it is from this that we get our words electron and electricity. What causes static electricity? Usually materials are electrically neutral. This means that they are "not" electrically charged. However this doesn't mean that materials have no charges "inside" them. All materials are made of atoms. Atoms are made of particles called electrons, protons, and neutrons. Electrons have a negative charge, protons have a positive charge, and neutrons (as you can probably guess by the name) have no charge. Materials are partly made of charged particles. Materials have no overall charge whenever they have the same number of electrons and protons. The negative charge of the electrons exactly cancels out the positive charge of the protons. Test YourSelf: Q1) The element lithium has three protons. How many electrons must a neutral lithium atom have? Inside an atom the protons and neutrons are held firmly at the very centre in a structure called the nucleus. The electrons however are held much more loosely. Some can be found at the surface, and if the surface is rubbed, they can be rubbed away. This removal of electrons will leave the atom with more protons than electrons. It will have a positive charge because the protons are positively charged. Atoms that are charged are called ions. When you rub a balloon on your sweater some of the electrons are transferred from the sweater to the balloon. The sweater becomes positively charged. Test YourSelf: Q2) What charge will the balloon have? How do the two types of charge behave? The sweater is positively charged, the balloon negative. The balloon sticks to the sweater. This is because unlike charges attract. If you were to take two balloons, both of which had been negatively charged and brought them together you would find that they tried to push apart. This is because like charges repel.  Test Yourself Q3) Look at the balloons above. The pair on the left are being repelled. The pair on the right are being attracted. Notice that one out of each pair of balloons is negatively charged. What is the sign (positive or negative) of the charges on the other two balloons? Sparks can fly! When a large enough charge is built up sparks can fly. A spark is just a flow of electrons through the air. This can sometimes be very dangerous. For example, petrol flowing through a pipeline can pick up a charge by friction. If a spark flies it could ignite the petrol. For this reason, pipes are earthed. Earthing means connecting the pipe to the ground by means of a copper wire. Copper, unlike air, is a very good conductor of electricity. The electrons can flow easily and safely to the ground without sparking. Some devices make use of sparks. A Van de Graaf generator consists of a large metal dome that is in contact with a rubber or plastic looped belt. As the belt moves, friction causes it to become charged; the charge is spread out over the metal dome where it builds up until it is large enough to cause a spark. In the diagram above a girl is standing on a washing up bowl and touching a Van de Graaf machine. Test Yourself Q4) What do you think is the purpose of the washing up bowl? Test Yourself Q5) Explain why the girl's hair is standing on end (no it's not her style, it's to do with the charge) Answers to questions | « Electricity | Uses of static electricity » 

Back to GCSE Science/Static Electricity Q1) 3 --- all atoms have the same number of protons as electrons Q2) Negative --- electrons have been added to the balloon. Q3 left hand pair) Negative--- The balloons are repelled Q3 right hand pair)Positive --- The balloons are attracted. Remember '"like charges repel, unlike charges attract Q4) The purpose of the washing up bowl is to insulate the girl from the ground. People are reasonable conductors of electricity. If the girl just stood on the floor electricity would run through her body and leak away. She would act as an earth wire for the Van de Graaf. Q5) The girl's hairs all pick up static electricity from the Van de Graaf machine. Remember like charges repel, so each hair tries to get as far away from the other hairs as possible. That is why her hair stands on end. 

GCSE Science/Electricity There are several practical uses of static electricity in our daily life. We will look at three of them on this page. There are many, many more, but these are the easiest to understand. Let’s take a look at them. Number One: The Photocopier. One example of the practical use of static electricity is a photocopier. A photocopier is a complicated piece of equipment, but the basic principle of how it works is fairly simple. The best way to understand what is going on is to consider it as a stage by stage process. Relates to xerography Stage one. Positive charge is applied onto a plate from a high voltage power supply which is called charging by friction. The plate is connected to the earth but the charge does not have quite enough energy to flow away from it. (The plate is not a good conductor of electricity.) Stage two. Paper is placed over the plate and a light shines onto the paper. Where the paper is white the light is reflected onto the plate. Where the paper is dark a shadow falls onto the plate. The light falling on the plate gives it just the extra energy needed to allow the charge to escape to earth. The plate becomes neutral where the paper is white but keeps its charge where the paper is black. The plate is now a copy of the paper with charges taking the place of ink. You could call this a template. Stage three. Toner particles are sprayed through a negatively charged nozzle onto the plate. As the toner passes through the nozzle it picks up the charge so that each particle of toner becomes negatively charged. The now charged toner is attracted to the areas of positive charge because unlike charges attract. More light then allows the positive charge to escape (However the negative charge on the toner remains.) Stage four. A sheet of paper is given a very strong positive charge, and then placed in contact with the plate. The paper attracts the toner. The paper is then removed from the plate and passed through a heating unit. The heat melts the toner and bonds it to the paper. In a real photocopier, there is no plate, just a large drum. As the drum rotates its surface goes through stages one through four. At the end of the sequence a scraper removes any toner left on the drum and the whole process is repeated with a new image. A good photocopier is capable of producing 20 duplicate pages per minute (20ppm), which is approximately one page every three seconds. Questions on Photocopiers. Q1) If the black areas of the image leave a positive charge in the plate what charge do the white areas of the image leave? ("Be very careful, the answer may not be what you think!") Q2) How does the light shining onto the charged plate allow it to lose its charge? Q3) Why is the toner given a negative charge? Q4) Why does the paper attract the toner? Other uses of static electricity. Number Two: Spray painting car parts. When paint is sprayed from a paint gun, the painter normally needs to use a fair amount of skill to ensure the paint goes on evenly. By connecting the spray nozzle to a negative electrode, it is possible to charge each droplet of paint. If the car part is then given the opposite charge, the paint droplets will be attracted to the car body part. This has several advantages: Questions. Q1) If the paint droplets are given a positive charge, what charge should the car part be given? A1) A negative charge as the positive will always be looking for a negative to balance it out. Q2) Why does less paint fall on the floor? Give reasons for your answer A2) This is because the oppositely charged particles are naturally attracted to each other and depending on the strength of the charge the range of attraction will increase/decrease respectively. Q3) In the application of painting cars, wouldn't negatively charging the spray source and positively charging the object being sprayed also attract dust, dirt, and trash too? Number Three: Pollution Control. Static electricity is used in pollution control by applying a static charge to dirt particles in the air and then collecting those charged particles on a plate or collector of the opposite electrical charge. Such devices are often called electrostatic precipitators. Out of the three this is the one people will know it most for. Factories use static electricity to reduce pollution coming from their smoke stacks. They give the smoke an electric charge. When it passes by electrodes of the opposite charge, most of the smoke particles cling to the electrodes. This keeps the pollution from going out into the atmosphere. Answers | «Static Electricity | Advanced topics» 

Introduction. A "graph" is a mathematical way of representing the concept of a "network". A network has points, connected by lines. In a graph, we have special names for these. We call these points "vertices" (sometimes also called nodes), and the lines, "edges". Here is an example graph. The edges are red, the vertices, black. In the graph, formula_1 are vertices, and formula_2 are edges. Definitions of graph. There are several roughly equivalent definitions of a graph. Set theory is frequently used to define graphs. Most commonly, a graph formula_3 is defined as an ordered pair formula_4, where formula_5 is called the graph's vertex-set and formula_6 is called the graph's edge-set. Given a graph formula_3, we often denote the vertex—set by formula_8 and the edge—set by formula_9. To visualize a graph as described above, we draw formula_10 dots corresponding to vertices formula_11. Then, for all formula_12 we draw a line between the dots corresponding to vertices formula_13 if and only if there exists an edge formula_14. Note that the placement of the dots is generally unimportant; many different pictures can represent the same graph. Alternately, using the graph above as a guide, we can define a graph as an ordered triple formula_15: In the above example, If formula_20 is not injective — that is, if formula_21 such that formula_22 — then we say that formula_3 is a multigraph and we call any such edges formula_24 "multiple edges". Further, we call edges formula_25 such that formula_26 loops. Graphs without multiple edges or loops are known as simple graphs. Graphs can, conceivably, be infinite as well, and thus we place no bounds on the sets V and E. We will not look at infinite graphs here. Directions, Weights, and Flows. We define a directed graph as a graph such that formula_20 maps into the set of ordered pairs formula_28 rather than into the family of two-element sets formula_29. We can think of an edge formula_17 such that formula_31 as 'pointing' from formula_32 to formula_33. As such we would say that formula_32 is the "tail" of edge formula_35 and that formula_33 is the "head". This is one of the vagaries of graph theory notation, though. We could just as easily think of formula_32 as the head and formula_33 as the tail. To represent a directed graph, we can draw a picture as described and shown above, but place arrows on every edge corresponding to its direction. In general, a weight on a graph formula_3 is some function formula_40. A flow formula_41 is a directed graph formula_15 paired with a weight function such that the weight "going into" any vertex is the same amount as the weight "going out" of that vertex. To make this more formal, define sets Then, formally stated, our requirement on the weight function is formula_45 Algebraic Graph Theory. While set theory is frequently used when discussing graphs, other approaches can simplify certain operations. A set can be defined using an adjacency matrix formula_46 where element formula_47 is a 1, if there is an edge between vertex i and vertex j and 0 otherwise. Special Graphs. Some graphs occur frequently enough in graph theory that they deserve special mention. One such graphs is the "complete graph" on n vertices, often denoted by Kn. This graph consists of n vertices, with each vertex connected to every other vertex, and every pair of vertices joined by exactly one edge. Another such graph is the "cycle graph" on "n" vertices, for "n" at least 3. This graph is denoted C"n" and defined by V := {1,2..,n} and E := . Even easier is the "null graph" on "n" vertices, denoted N"n"; it has "n" vertices and no edges! Note that N1 = K1 and C3 = K3. Some Terms. Two vertices are said to be "adjacent" if there is an edge joining them. The word "incident" has two meanings: Two graphs "G" and "H" are said to be "isomorphic" if there is a one-to-one function from (or, if you prefer, one-to-one correspondence between) the vertex set of "G" to the vertex set of "H" such that two vertices in "G" are adjacent if and only if their images in "H" are adjacent. (Technically, the multiplicity of the edges must also be preserved, but our definition suffices for simple graphs.) Subgraphs. A "subgraph" is a concept akin to the subset. A subgraph has a subset of the vertex set V, a subset of the edge set E, and each edge's endpoints in the larger graph has the same edges in the subgraph. A A subgraph formula_48 of formula_3 is "generated" by the vertices {formula_50}formula_51 if the edge set of formula_48 consists of all edges in the edge set of formula_3 that joins the vertices in formula_54{formula_55}. A "path" is a sequence of edges formula_56 such that ei is adjacent to ei+1 for all i from 1 to N-1. Two vertices are said to be connected if there is a path connecting them. Trees and Bipartite Graphs. A "tree" is a graph that is (i) connected, and (ii) has no cycles. Equivalently, a tree is a connected graph with exactly formula_57 edges, where there are formula_10 nodes in the tree. A "Bipartite graph" is a graph whose nodes can be partitioned into two disjoint sets U and W such that every edge in the graph is incident to one node in U and one node in W. A tree is a bipartite graph. A "complete bipartite graph" is a bipartite graph in which each node in U is connected to every node in W; a complete bipartite graph in which U has formula_10 vertices and V has formula_60 vertices is denoted formula_61. Adjacent,Incident,End Vertices Self loops,Parallel edges,Degree of Vertex Pendant Vertex : Vertex Degree one "Pendant Vertex" Isolated Vertex : Vertex Degree zero "Isolated Vertex" Hamiltonian and Eulerian Paths. Hamiltonian Cycles: A Hamiltonian Cycle received its name from Sir William Hamilton who first studied the travelling salesman problem. A Hamiltonian cycle is a path that visits every vertex once and only once i.e. it is a walk, in which no edge is repeated (a trail) and therefore a trail in which no vertex is repeated (a path). Note also it is a cycle, the last vertex is joined to the first. A graph is said to be Eulerian if it is possible to traverse each edge once and only once, i.e. it has no odd vertices or it has an even number of odd vertices (semi-Eulerian). This has implications for the Königsberg problem. It may be easier to imagine this as if it is possible to trace the edges of a graph with a pencil without lifting the pencil off the paper or going over any lines. Planar Graphs. A "planar graph" is an undirected graph that can be drawn on the plane or on a sphere in such a way that no two edges cross, where an edge formula_62 is drawn as a continuous curve (it need not be a straight line) from u to v. Kuratowski proved a remarkable fact about planar graphs: A graph is planar if and only if it does not contain a subgraph homeomorphic to formula_63 or to formula_64. (Two graphs are said to be homeomorphic if we can shrink some components of each into single nodes and end up with identical graphs. Informally, this means that non-planar-ness is caused by only two things—namely, having the structure of formula_63 or formula_64 within the graph). Coloring Graphs. A graph is said to be planar if it can be drawn on a plane in such way that no edges cross one another except of course for meeting at vertices Each term, the Schedules Office in some university must assign a time slot for each final exam. This is not easy, because some students are taking several classes with finals, and a student can take only one test during a particular time slot. The Schedules Office wants to avoid all conflicts, but to make the exam period as short as possible. We can recast this scheduling problem as a question about coloring the vertices of a graph. Create a vertex for each course with a final exam. Put an edge between two vertices if some student is taking both courses. For example, the scheduling graph might look like this: Next, identify each time slot with a color. For example, Monday morning is red, Monday afternoon is blue, Tuesday morning is green, etc. Assigning an exam to a time slot is now equivalent to coloring the corresponding vertex. The main constraint is that adjacent vertices must get different colors; otherwise, some student has two exams at the same time. Furthermore, in order to keep the exam period short, we should try to color all the vertices using as few different colors as possible. For our example graph, three colors suffice: red, green, blue. The coloring corresponds to giving one final on Monday morning (red), two Monday afternoon (blue), and two Tuesday morning (green)... K Coloring. Many other resource allocation problems boil down to coloring some graph. In general, a graph G is kcolorable if each vertex can be assigned one of k colors so that adjacent vertices get different colors. The smallest sufficient number of colors is called the chromatic number of G. The chromatic number of a graph is generally difficult to compute, but the following theorem provides an upper bound: Theorem 1. A graph with maximum degree at most k is (k + 1)colorable. Proof. We use induction on the number of vertices in the graph, which we denote by n. Let P(n) be the proposition that an nvertex graph with maximum degree at most k is (k + 1)colorable. A 1 vertex graph has maximum degree 0 and is 1colorable, so P(1) is true. Now assume that P(n) is true, and let G be an (n + 1)vertex graph with maximum degree at most k. Remove a vertex v, leaving an nvertex graph G . The maximum degree of G is at most k, and so G is (k + 1)colorable by our assumption P(n). Now add back vertex v. We can assign v a color different from all adjacent vertices, since v has degree at most k and k + 1 colors are available. Therefore, G is (k + 1)colorable. The theorem follows by induction. Weighted Graphs. A weighted graph associates a label (weight) with every edge in the graph. Weights are usually real numbers, and often represent a "cost" associated with the edge, either in terms of the entity that is being modeled, or an optimization problem that is being solved. 

About the platform. The GNU operating system was started by Richard Stallman as a free replacement for the UNIX operating system. At the same time Linus Torvalds was working on a kernel, which he adapted to fit the GNU operating system. As time progressed, many applications from UNIX and DOS were ported to GNU/Linux as well as the thousands of new applications written for it. GNU/Linux has become a completely self-sufficient operating system with applications ranging from the many console applications to the numerous highly advanced GUI applications (many of which are based on lower level console applications) and everything else in between. The most popular languages for use on the GNU/Linux platform include C/C++ and , however the range of programming languages supported by the GNU/Linux platform cover the entire spectrum of the software development world. Other popular languages are Perl, Python and Ruby. Shell scripting is often used for administrative tasks but cannot be called a complete high level language. Basic information. Most UNIX code is instantly portable to GNU/Linux systems - it can be compiled as on a UNIX system. GNU/Linux programming tools are mostly from the GNU project at http://www.gnu.org, including gcc (free C/C++ compiler), and equivalents of make, ld, as, etc. Numerous others are available for various languages, including java. 

^Lesson 4^ Fill in the blank. Please fill in the blank with the correct form of the verb: 1. Yo _____________________ (cerrar) la puerta. "I close the door." 2. Tú _______________________ (perder) tus llaves. "You lost your keys." 3. ¿_________________________(preferir) usted la langosta? "Do you prefer lobster?" 4. Nosotros _______________________ (entender) los ejercicios. "We understand the exercises." 5. Felipe y Juan ______________________(empezar) las prácticas de gramática. "Felipe and Juan begin the grammar practices." 6. Ellos ____________________ (querer) almorzar. "They want to eat lunch." 7. La chica _____________________ (perder) su pañuelo. "The girl loses her handkerchief." 8. ¿Qué ______________________ (pensar)? "What do you -singular- think?" 9. Sofía _______________________ (entender) a Julio. "Sofía understands Julio." 10. Yo no ______________________ (querer) la ensalada. "I don't want the salad" Bonus: Argentinian voseo 2. Vos _____________________ (perder) tus llaves. "You lose your keys." 8. ¿Qué ______________________ (pensar)? "What do you -singular- think?" Soluciones a los ejercicios "Solutions to the exercises" ^Lesson 4^ 

« Catalysis | pKa values » Metabolism. Anabolism and catabolism. Metabolism (Fig. 1) is, broadly speaking, the conversion of food into energy, cell components, and waste products. &lt;br&gt; Figure 1: Overview of metabolism The above diagram shows the different parts of metabolism: Catabolic reactions release energy and are therefore exergonic, while anabolic reactions use up energy and are therefore endergonic. High-energy phosphates. Due to the large variety of food compounds, and the large number of biochemical reactions which need energy in anabolism, it would be quite inefficient to couple a specific anabolic reaction to a specific energy source in catabolism. Instead, the cell uses an intermediate compound, a kind of universal energy currency. This intermediate is called "high-energy phosphate". But when is a phosphate group called "high-energy", and how does it differ from a "low-energy" phosphate? A giveaway is the ΔG0' of hydrolysis. Hydrolysis separates a phosphate from a compound by adding water:     O                      O  R-OP-OH + H2O ⇌ R-OH + HO-P-OH     O                      O The ΔG0' of a low-energy (or "inorganic") phosphate group (called Pi) is 9-20 kJ mol-1, while the ΔG0' of a high-energy phosphate (denoted Ⓟ) is ~30 kJ mol-1. pKa value. Now what makes this Ⓟ so special? To explain this, we must take a little excourse into pH and pKa values. A phosphate group has between zero and three OH groups. This allows Ⓟ to exist in up to four different forms (0, 1, 2, and 3 OH groups, Fig. 2), depending on the pH value of the surrounding solution. A pKa value gives us the pH value at which 50% of the molecules are in one form (e.g., 1 OH group) and another (e.g., 2 OH groups). This is expressed by the "Henderson-Hasselbalch equation" : &lt;br&gt; Figure 2: The four possible forms of a phosphate group. pKa2 represents the conditions in the cell. Now to the promised difference between Ⓟ and PPi. The breaking of the ester bond of an ROⓅ releases more energy than the breaking of a PPi bond (Fig. 3), because of &lt;br&gt; Figure 3: Hydrolysis of Ⓟ and PPi. &lt;br&gt; Figure 4: Resonance stabilization of Pi. Resonance stabilization means that both OH and =O can "travel" around the phosphate. Of course, this is a crude analogy; they do not really move, the electrons are just "smeared" around the phosphate atom. This is also indicated by the use of the ↔ arrow, instead of ⇌; the three forms do not exist, they are just a way of writing down the chemical reality. As you can see in Fig. 3, the ΔG0' value for PPi⇌2Pi is ≪0, shifting the reaction strongly in favor of the 2Pi. Molecules using high-energy phosphates. Anhydride between phosphoric acid and carboxyl group. Hydrolysis : ΔG0' = -49.3 kJ mol-1&lt;br&gt; Guanidine phosphate. Hydrolysis : ΔG0' = -43.0 kJ mol-1&lt;br&gt; Enol phosphate. In the below picture, the final product should not have a carbon-carbon double bond, but a single bond with CH3 on the top. It is an error. Hydrolysis : ΔG0' = -61.9 kJ mol-1&lt;br&gt; ATP. Adenosine triphosphate contains one low-energy and two high-energy phosphate bonds:&lt;br&gt; &lt;br&gt; Low energy : ΔG0' = -14,2 kJ mol-1&lt;br&gt; High energy : ΔG0' = -30.5 kJ mol-1&lt;br&gt; Basically, any ATP-driven reaction is reversible, building ATP from ADP and Pi in the process. However, some ATP-driven reactions should never be reversed; these include nucleotide and protein synthesis. If these were reversed, the organism would disassemble its own DNA and proteins for energy, a rather unfortunate strategy. For reactions that should never be reversed, ATP can be broken down into AMP (adenosine monophosphate) and PPi, which in turn becomes 2×Pi. This reaction has a ΔG0' of -65,7 kJ mol-1, which is totally irreversible under "in vivo" conditions. It should be noted that AMP can not directly be converted to ATP again. Instead, the enzyme "AMP kinase" forms two ADP molecules from one ATP and one AMP. The resulting ADPs are then treated as described above. Non-covalent bonds. The destruction of covalent bonds takes up huge amounts of energy. The breakdown of an O2 molecule into two oxygen atoms needs ~460 kJ mol-1. Thus, nowhere in "living" biochemistry are covalent bonds actually destroyed; if one is broken, another one is created. This is where non-covalent bonds come in, they are weak enough to be broken down easily, and to form "bonds" again. For this reason, many biochemical functions are using so-called weak/secondary/non-covalent bonds. Weak bonds are created and destroyed much more easily than covalent ones. The typical range of energy needed to destroy such a weak bond is 4-30 kJ mol-1. Thus, the formation of weak bonds is energetically favorable, but these bonds are also easily broken by kinetic (thermal) energy (the normal movement of molecules). Biochemical interactions are often temporary (e.g., a substrate has to leave an enzyme quickly after being processed), for which the weakness of these bonds is essential. Also, biochemical specificity (e.g., enzyme-substrate-recognition) is achieved through weak bonds, utilizing two of their major properties: The link that follows demonstrates the type of non-covalent forces: There are three basic types of weak bonds, and a fourth "pseudo-bond": Ionic bonds. Ionic bonds are electrostatic attractions between permanently charged groups. Ionic bonds are not directed. Example: Hydrogen bonds. Hydrogen bonds are also established by electrostatic attraction. These attractions do not occur between permanently charged groups, but rather between atoms temporarily charged by a "dipole moment", resulting from the different electronegativity of atoms within a group. Hydrogen bonds are even weaker than ionic bonds, and they are highly directional, usually along a straight line. Besides being weaker than ionic bonds, hydrogen bonds are also weaker, and longer than similar covalent bonds. Hydrogen bonds are unique because they only exist when the Hydrogen is bonded to an oxygen (O), Nitrogen (N), or Fluorine (F), but the most common hydrogen bonds in biochemistry are: Hydrogen bonds equal an energy between 12-29 kJ mol, whereas covalent bonds are much higher. For example, the covalent bond between oxygen and hydrogen is about 492 KJ mol-1. Hydrogen Bonds and Water. Water has unique properties; after all, it is chosen to be the universal solvent. The unique properties of water are due to hydrogen bonding between all the oxygen and hydrogen atoms of the content. The hydrogen bonds occurring in water are about 2 angstroms apart from each other. Although hydrogen bonding is only about 5% as strong as covalent bond, they still cause water to have a high boiling point, and a high surface tension. The following link will take you to the structure of water and its Hydrogen Bonding. Van der Waals attractions. Van der Waals attractions are established between electron density-induced dipoles. They form when the outer electron shells of two atoms "almost" (but not quite) touch. The distance of the atoms is very important for these weak interactions. If the atoms are too far apart, the interactions are too weak to establish; if the atoms are too close to each other, their electron shells will repel each other. Van der Waals attractions are highly unspecific; they can occur between virtually any two atoms. Their energy is between 4-8 kJ mol-1. Hydrophobic interactions. Hydrophobic forces are not actually bonds, so this list has four items, but still just three bond types. In a way hydrophobic forces are the negation of the hydrogen bonds of a polar solute, usually water, enclosing a nonpolar molecule. For a polar solute like water, it is energetically unfavorable to "waste" a possible hydrogen bond by exposing it towards a nonpolar molecule. Thus, water will arrange itself around any nonpolar molecule in such a way that no hydrogen bonds point towards that molecule. This results in a higher order, compared to "freely" moving water, which leads to a lower entropy level and is thus energetically unfavorable. If there is more than one nonpolar molecule in the solute, it is favorable for the nonpolar molecules to aggregate in one place, reducing their surrounding, ordered "shell" of water to a minimal surface. Also, in large molecules, such as proteins, the hydrophobic (nonpolar) parts of the molecule will tend to turn towards the inside, while the polar parts will tend to turn towards the surface of the molecule. References. Cooke, Rosa-lee. "Properties of Water". Lecture 10. Mountain Empire Community College. n.d. Web. http://water.me.vccs.edu/courses/env211/lesson10_print.htm Kimball, John W.. "Hydrogen Bonds". Kimball’s Biology Pages. Feb. 12, 2011. Web. http://users.rcn.com/jkimball.ma.ultranet/BiologyPages/H/Hbonds_water.gif Lower, Stephen. "States of matter: Water and hydrogen bonding". General Chemistry Virtual Textbooks. 2009. Aug. 26, 2010. Web. http://www.chem1.com/acad/webtext/states/water.html n.p. "Covalent vs. Non-Covalent Bonds". n.d. http://www.pearsonhighered.com/mathews/ch02/c02cv.htm W. W. Norton &amp; Company. "Hydrogen Bonding in Water". Web. 2012. http://www.wwnorton.com/college/chemistry/gilbert2/tutorials/chapter_10/water_h_bond/ WyzAnt Tutoring. "WyzAnt Tutoring". Bonds. 2012. Web. http://www.wyzant.com/Help/Science/Chemistry/Bonds/ « Catalysis | pKa values » 

Cell Biology | Parts of the cell | Organelles « Golgi apparatus | Lysosomes | Peroxisomes » Lysosomes. Membrane-bound sacs called lysosomes contain digestive enzymes that can break down such macromolecules as proteins, nucleic acids, and polysaccharides (Figure 6-23). Lysosomes have several functions. They fuse with incoming food vacuoles and expose the nutrients to enzymes that digest them, thereby nourishing the cell. Lysosomes also function like safety officers when they help destroy harmful bacteria. In certain cells—for example, your white blood cells—lysosomes release enzymes into vacuoles that contain trapped bacteria and break down the bacterial cell walls. Similarly, lysosomes serve as recycling centers for damaged organelles. Without harming the cell, a lysosome can engulf and digest another organelle. This makes molecules available for the construction of new organelles.The structures vary in size from 0.2 to 2 micrometers in diameter. The staining reveals a crystal like matrix in spherical vesicles. The crystalloid matrix is urate oxidase. These are small organelles containing around 40 enzymes for intercellular digestion. The lysosome membrane helps to protect the enzymes as much as it helps protect the cell. This is because the optimal pH for these enzymes is around a pH of 5. The membrane of the lysosome is again a lipid bilayer and is thought to have a ATP hydrolysis to pump H+ into the lysosome to maintain the pH. This also has another affect, that is free protons. Other small molecules can pass through the lysosome membrane, but will then become charged by picking up a free proton, then they are less likely to be able to leave the lysososome. A good reference on Lysosomes is at « Golgi apparatus | Lysosomes | Peroxisomes » Cell Biology | Parts of the cell Organelles 

Cell Biology | ../Parts of the cell/ | ../Organelles/ « Lysosomes | Peroxisomes | Cytosol » It used to be thought that peroxisomes are formed by the budding of smooth Endoplasmic Reticulum (ER). However, now it is thought that they form through self-assembly. (I will get more information references as I do some more thorough literature searches). The peroxisome is another major source of Oxygen utilization (along with the mitochondrion). There are specific proteins associated with the peroxisomes membrane, also there are 3 oxidation enzymes associated with peroxisomes: The enzyme contents vary with various types of cells. One of the main functions of peroxisomes in liver cells is detoxification. This is done by the oxidation of substances like: Why peroxisomes are not like lysosomes. Peroxisomes are organelles that contain oxidative enzymes, such as D-amino acid oxidase, urate oxidase, and catalase. They may resemble a lysosome, however, they are not formed in the Golgi complex. Peroxisomes are distinguished by a crystalline structure inside a sac which also contains amorphous gray material. They are self replicating, like the mitochondria. Components accumulate at a given site and they can be assembled into a peroxisome. They may look like storage granules, however, they are not formed in the same way as storage granules. Peroxisomes function to rid the body of toxic substances like hydrogen peroxide, or other metabolites. They are a major site of oxygen utilization and are numerous in the liver where toxic byproducts are going to accumulate. The peroxisome is made as a phospholipid bilayer, encapsulating oxidative materials. They would be 'sphere-ish' in shape, not necessarily a perfect sphere, and sometimes, they may take other shapes. But most electron micrographs I have seen (2 dimensions) show them as circles. (As you may be aware, the Cell membrane is also a phospholipid bilayer.) Peroxisomes have membrane proteins that are critical for peroxisomal function, to import proteins into their interiors, proliferate or segregate to daughter cells (This update, thanks to Babich Temps). The main differences would be: « Lysosomes | Peroxisomes | Cytosol » Cell Biology | ../Parts of the cell/ ../Organelles/ 

Cell Biology | ../Parts of the cell/ | ../Organelles/ « Peroxisomes | Cytosol | Cytoskeleton and Microtubules » The cytosol (as opposed to cytoplasm, which also includes the organelles) is the internal fluid of the cell, and a large part of cell metabolism occurs here. Proteins within the cytosol play an important role in signal transduction pathways, glycolysis, and act as intracellular receptors and ribosomes. In prokaryotes, all chemical reactions take place in the cytosol. In eukaryotes, the cytosol contains the cell organelles. In plants, the amount of cytosol can be reduced due to the large tonoplast (central vacuole) that takes up most of the room of the cell. The cytosol is not a "soup" with free-floating particles, but highly organized on the molecular level. The cytosol also contains the cytoskeleton. It is made of fibrous proteins and (in many organisms) maintains the shape of the cell, anchors organelles, and controls internal movement of structures, e.g., transport vesicles. As the concentration of soluble molecules increases within the cytosol, an osmotic gradient builds up toward the outside of the cell. Water is flowing into the cell, making it larger. To prevent the cell from bursting apart, molecular pumps in the plasma membrane, the cytoskeleton, the tonoplast or a cell wall (if present) are used to counteract the osmotic pressure. « Peroxisomes | Cytosol | Cytoskeleton and Microtubules » Cell Biology | ../Parts of the cell/ | ../Organelles/ 

GCSE Science/Electricity Before reading this module check to see if you need to. If you intend to take the foundation paper, you may find that you do not need to do the work on this page. If in doubt check with your teacher. Where charge concentrates. When an object is charged up with electrons, the electrons try to spread out over the object to be as far apart as possible. This means they go to the surface rather than spread throughout. On a metal object the material of which the object is made shields the electrons from one another. They do not ‘’see’’ the electrons round the back. This means that highly curved objects can hold more electrons than flat objects. For a complicated shape, the electrons tend to congregate on the more highly curved areas, and to desert the flatter areas. Inducing charge separation on neutral objects. Consider a crystal of gold for example. {This argument works for all materials}. Normally the atoms are perfectly spherical, and completely neutral. Q1) Why are the atoms neutral? Remember that neutral atoms contain positive charges in the nucleus and negative electrons orbiting in their shell. Normally the center of the positive charges and the center of the negative electrons is in the same place: the exact center of the atom. If we bring up a charged object to them however, something important happens. Look at the diagram on the right. A charged rod is brought near the surface of the crystal. {Note: the rod is not shown, just the charges on it}. The electrons in the atoms try to get away from the negative charge because they are repelled. The positive charges in the nucleus try to get nearer to the negative rod because they are attracted to it. The result is the atom slightly changes shape. The center of the negative charges goes to the right of the center of the atom and the center of the positive charges goes to the left. We say there is a separation of charges Q2) What would happen if we brought up a positively charged rod? There is always an attractive force between the rod and the crystal no matter what the charge on the rod. To see why this is so, we need to know one more thing about electrostatics. The force of repulsion or attraction falls off with distance. In fact, rather like gravity it falls off as the "square" of the distance. This means if you go to twice the distance the force becomes only a quarter as much. If you go three times the distance it becomes only one ninth as great and so on. Now look at the diagram again. The negative charges in the atoms are repelled by the rod, the positive ones are attracted, but the positive charges are closer, so overall there is a force of attraction. Q3) Repeat the above argument with a positively charged rod to convince yourself that there is still an attractive force. Lightning. Lightning is basically a great big spark. In a thunderstorm, separation of charge occurs in big rain clouds. No one is really sure of "how" this occurs, although it appears to be caused by ice crystals rubbing together and becoming charged. but what "is" known is that storm clouds have a negative bottom and a positive top. The bottom of the cloud is nearer to the ground so it induces a positive charge on the ground. When the voltage becomes high enough a spark flies between the cloud and the earth; this is lightning. The spark is so hot, it causes the air to rapidly expand then collapse; this causes thunder. Q4) Why does the voltage need to be very high before lightning can occur? Lightning rods. Lightning rods are used to protect buildings from lightning strikes. They are usually made from a length of copper rod with a pointed end. They are attached to tall buildings and the lightning strikes them rather than the building itself. Q5) Why do you think the rods are made of copper rather than iron which is cheaper? Q6) Why do the rods have a pointed end? Q7) It is standard advice that if caught out in a storm you should "never" stand under a lone tree. Why is this advice good? Answers | «Uses of static electricity | Electrolysis» 

Welcome to the Wikibook of GEOMETRY Preface. The word geometry originates from the Greek words ("geo" meaning world, "metri" meaning measure) and means, literally, to measure the earth. It is an ancient branch of mathematics, but its modern meaning depends largely on context. Geometry largely encompasses forms of non-numeric mathematics, such as those involving measurement, area and perimeter calculation, and work involving angles and position. It was one of the two fields of pre-modern mathematics, the other being the study of numbers. In modern times, geometric concepts have been generalized to a high level of abstraction and complexity, and have been subjected to the methods of calculus and abstract algebra, so that many modern branches of the field are barely recognizable as the descendants of early geometry. This Wikibook is dedicated to high school geometry and geometry in general. High School Geometry. The outline of topics reflects the California curriculum content standards.  

In C, you have already considered creating variables for use in the program. You have created some arrays for use, but you may have already noticed some limitations: "Dynamic memory allocation" in C is a way of circumventing these problems. The codice_1 function. void *calloc(size_t nmemb, size_t size); void free(void *ptr); void *malloc(size_t size); void *realloc(void *ptr, size_t size); The standard C function codice_1 is the means of implementing dynamic memory allocation. It is defined in stdlib.h or malloc.h, depending on what operating system you may be using. Malloc.h contains only the definitions for the memory allocation functions and not the rest of the other functions defined in stdlib.h. Usually you will not need to be so specific in your program, and if both are supported, you should use &lt;stdlib.h&gt;, since that is ANSI C, and what we will use here. The corresponding call to release allocated memory back to the operating system is codice_3. When dynamically allocated memory is no longer needed, codice_3 should be called to release it back to the memory pool. Overwriting a pointer that points to dynamically allocated memory can result in that data becoming inaccessible. If this happens frequently, eventually the operating system will no longer be able to allocate more memory for the process. Once the process exits, the operating system is able to free all dynamically allocated memory associated with the process. Let's look at how dynamic memory allocation can be used for arrays. Normally when we wish to create an array we use a declaration such as int array[10]; Recall codice_5 can be considered a pointer which we use as an array. We specify the length of this array is 10 codice_6s. After codice_7, nine other integers have space to be stored consecutively. Sometimes it is not known at the time the program is written how much memory will be needed for some data; for example, when it depends upon user input. In this case we would want to dynamically allocate required memory after the program has started executing. To do this we only need to declare a pointer, and invoke codice_1 when we wish to make space for the elements in our array, "or", we can tell codice_1 to make space when we first initialize the array. Either way is acceptable and useful. We also need to know how much an int takes up in memory in order to make room for it; fortunately this is not difficult, we can use C's builtin codice_10 operator. For example, if codice_11 yields 4, then one codice_6 takes up 4 bytes. Naturally, codice_13 is how much memory we need for 2 codice_6s, and so on. So how do we codice_1 an array of ten codice_6s like before? If we wish to declare and make room in one hit, we can simply say int *array = malloc(10*sizeof(int)); We only need to declare the pointer; codice_1 gives us some space to store the 10 codice_6s, and returns the pointer to the first element, which is assigned to that pointer. Important note! codice_1 does "not" initialize the array; this means that the array may contain random or unexpected values! Like creating arrays without dynamic allocation, the programmer must initialize the array with sensible values before using it. Make sure you do so, too. ("See later the function codice_20 for a simple method.) It is not necessary to immediately call codice_1 after declaring a pointer for the allocated memory. Often a number of statements exist between the declaration and the call to codice_1, as follows: int *array = NULL; printf("Hello World!!!"); /* more statements */ array = malloc(10*sizeof(int)); /* delayed allocation */ /* use the array */ A more practical example of dynamic memory allocation would be the following:Given an array of 10 integers, remove all duplicate elements from the array, and create a new array without duplicate elements (a ).A simple algorithm to remove duplicate elements:  int arrl = 10; // Length of the initial array  int arr[10] = {1, 2, 2, 3, 4, 4, 5, 6, 5, 7}; // A sample array, containing several duplicate elements  for (int x = 0; x &lt; arrl; x++)  for (int y = x + 1; y &lt; arrl; y++)  if (arr[x] == arr[y])  for (int s = y; s &lt; arrl; s++)  if (!(s + 1 == arrl))  arr[s] = arr[s + 1];  arrl--;  y--; Because the length of our new array depends on the input, it must be dynamically allocated: int *newArray = malloc(arrl*sizeof(int)); The above array will currently contain unexpected values, so we must use codice_23 to set our dynamically allocated memory block to the new values: memcpy(newArray, arr, arrl*sizeof(int)); Error checking. When we want to use codice_1, we have to be mindful that the pool of memory available to the programmer is "finite". Even if a modern PC will have at least an entire gigabyte of memory, it is still possible and conceivable to run out of it! In this case, codice_1 will return codice_26. In order to stop the program crashing from having no more memory to use, one should always check that malloc has not returned codice_26 before attempting to use the memory; we can do this by int *pt = malloc(3 * sizeof(int)); if(pt == NULL)  fprintf(stderr, "Out of memory, exiting\n");  exit(1); Of course, suddenly quitting as in the above example is not always appropriate, and depends on the problem you are trying to solve and the architecture you are programming for. For example, if the program is a small, non critical application that's running on a desktop quitting may be appropriate. However if the program is some type of editor running on a desktop, you may want to give the operator the option of saving their tediously entered information instead of just exiting the program. A memory allocation failure in an embedded processor, such as might be in a washing machine, could cause an automatic reset of the machine. For this reason, many embedded systems designers avoid dynamic memory allocation altogether. The codice_28 function. The codice_28 function allocates space for an array of items and initializes the memory to zeros. The call codice_30 allocates codice_31 objects, each of whose size is sufficient to contain an instance of the structure codice_32. The space is initialized to all bits zero. The function returns either a pointer to the allocated memory or, if the allocation fails, codice_26. The codice_34 function.  void * realloc ( void * ptr, size_t size ); The codice_34 function changes the size of the object pointed to by codice_36 to the size specified by codice_37. The contents of the object shall be unchanged up to the lesser of the new and old sizes. If the new size is larger, the value of the newly allocated portion of the object is indeterminate. If codice_36 is a null pointer, the codice_34 function behaves like the codice_1 function for the specified size. Otherwise, if codice_36 does not match a pointer earlier returned by the codice_28, codice_1, or codice_34 function, or if the space has been deallocated by a call to the codice_3 or codice_34 function, the behavior is undefined. If the space cannot be allocated, the object pointed to by codice_36 is unchanged. If codice_37 is zero and codice_36 is not a null pointer, the object pointed to is freed. The codice_34 function returns either a null pointer or a pointer to the possibly moved allocated object. The codice_3 function. Memory that has been allocated using codice_1, codice_34, or codice_28 must be released back to the system memory pool once it is no longer needed. This is done to avoid perpetually allocating more and more memory, which could result in an eventual memory allocation failure. Memory that is not released with codice_3 is however released when the current program terminates on most operating systems. Calls to codice_3 are as in the following example. int *myStuff = malloc( 20 * sizeof(int)); if (myStuff != NULL)  /* more statements here */  /* time to release myStuff */  free( myStuff ); free with recursive data structures. It should be noted that codice_3 is neither intelligent nor recursive. The following code that depends on the recursive application of free to the internal variables of a struct does not work. typedef struct BSTNode  int value;  struct BSTNode* left;  struct BSTNode* right; } BSTNode; // Later: ... BSTNode* temp = (BSTNode*) calloc(1, sizeof(BSTNode)); temp-&gt;left = (BSTNode*) calloc(1, sizeof(BSTNode)); // Later: ... free(temp); // WRONG! don't do this! The statement "codice_58" will not free codice_59, causing a memory leak. The correct way is to define a function that frees "every" node in the data structure: void BSTFree(BSTNode* node){  if (node != NULL) {  BSTFree(node-&gt;left);  BSTFree(node-&gt;right);  free(node); Because C does not have a garbage collector, C programmers are responsible for making sure there is a codice_60 exactly once for each time there is a codice_61. If a tree has been allocated one node at a time, then it needs to be freed one node at a time. Don't free undefined pointers. Furthermore, using codice_3 when the pointer in question was never allocated in the first place often crashes or leads to mysterious bugs further along. To avoid this problem, always initialize pointers when they are declared. Either use codice_1 at the point they are declared (as in most examples in this chapter), or set them to codice_26 when they are declared (as in the "delayed allocation" example in this chapter). Write constructor/destructor functions. One way to get memory initialization and destruction right is to imitate object-oriented programming. In this paradigm, objects are constructed after raw memory is allocated for them, live their lives, and when it is time for them to be destructed, a special function called a destructor destroys the object's innards before the object itself is destroyed. For example: /* this is the type of object we have, with a single int member */ typedef struct WIDGET_T {  int member; } WIDGET_T; /* functions that deal with WIDGET_T */ /* constructor function */ void WIDGETctor (WIDGET_T *this, int x)  this-&gt;member = x; /* destructor function */ void WIDGETdtor (WIDGET_T *this)  /* In this case, I really don't have to do anything, but  if WIDGET_T had internal pointers, the objects they point to  would be destroyed here. */  this-&gt;member = 0; /* create function - this function returns a new WIDGET_T */ WIDGET_T * WIDGETcreate (int m)  WIDGET_T *x = 0;  x = malloc (sizeof (WIDGET_T));  if (x == 0)  abort (); /* no memory */  WIDGETctor (x, m);  return x; /* destroy function - calls the destructor, then frees the object */ void WIDGETdestroy (WIDGET_T *this)  WIDGETdtor (this);  free (this); /* END OF CODE */ 

A "lattice" is a "poset" such that each pair of elements has a unique "least upper bound" and a unique "greatest lower bound". 

[[GCSE Science/Electricity]] Electrolysis is the decomposition of certain types of substance using electricity. The types of substance that can be split are ionic substances. This just means that they are made of charged ions rather than neutral atoms. {Remember that an ion is just an atom that has either a positive or negative charge}. An example of an ionic substance is common table salt sodium chloride. The sodium atom has a positive charge, the chlorine atom has a negative charge. It is usually written as Na+Cl-. Q1) Check in a periodic table, what is the symbol for sodium: Na or Cl? As you may already know if you've studied the Metals module, a salt is any substance made by combining an acid with an alkali. Acids, alkalis, and therefore all salts are ionic. Q2) Which of the following substances can be broken up by electricity: sodium chloride, iron sulphate, copper nitrate? Basic experimental setup. Most ionic compounds are not liquid at room temperature. This is a problem because the ions need to be able to move for the electric current to be able to flow. This can be achieved by melting. Look at the electrical setup shown on the right. The electrodes are just two carbon rods connected to a battery. The one connected to the positive electrode is called the anode. The one connected to the negative electrode is called the cathode. this is due to a collision Consider for example the compound lead bromide. This compound is a solid at room temperature but can be molten over a Bunsen flame. So what you would do is put some lead bromide into a beaker. Put the beaker on a tripod over a bunsen flame. Melt the lead bromide, then put in the electrodes and turn the power supply on at a setting of, say, 2V. What you would see happening is the cathode, is a silvery coating of pure lead forming, and bromine forming at the anode. The current would continue to flow until all the lead bromide was turned into lead and bromine. Q3) It takes energy to split up a compound like lead bromide. Where does this energy come from? Q4) Predict what products you would get at the anode and cathode if copper chloride was the electrolyte. What happens at the anode. The anode is the positive electrode; it attracts negatively charged ions, because unlike charges attract. The bromine ions move through the melt until they reach the anode. Once they get there, they give up their two extra electrons to become bromine atoms. 2Br- → Br2 + 2e- The electrons flow up the anode to the positive terminal of the battery. What happens at the cathode. The cathode is the negative electrode; it attracts the positively charged ions. Metal ions are always positive and so the lead ions flow through the metal uhe negatively charged terminal of the battery and onto the lead ions. Some trick to remember cations and anions, cathodes and anodes. I have a cat...I call her by saying come here plussy! - cathodes attract positive ions ca+ions has a plus in it, cations are positive ions red cat: reduction occurs at the cathode Pb2+ + 2e- → Pb Q5) Solid ionic substances do not conduct electricity and are not split up by it. Why do you think that is? Quantity calculations (higher tier only). In the experiment with lead bromide, you saw that lead was deposited at the cathode. If you actually do the experiment you will see that the lead coats the cathode. In this section we will look at how much metal will coat a cathode in a given time. A scientist performed the following experiment. His results were: You can see from the results that the total amount of copper deposited depends on both the current and the time it flows. This is because the number of copper atoms that can be made from ions depends on the total amount of charge that flows. The unit of charge is the coulomb. One coulomb is the amount of charge when one Ampere flows for one second. Q6) Look at the results table above. How much copper is deposited when 1A flows for 3000 seconds? Q7) How much copper do you predict would be deposited if 1A were to flow for 6000 seconds. Q8) What about if 2A were to flow for 12000 seconds ? Electrolysis of Aqueous solutions (Advanced). "Before studying this section check with your teacher to see if you need to". Earlier on in this module you've learned that ions must be able to move in order for electrolysis to work. If the ions are held rigid {such as in a solid}, they can't move and no electricity will flow. We've looked at how the freeing up of ions can occur by melting the electrolyte. Another way to achieve this is by dissolving the electrolyte in water. The trouble with this method is, there will be more than one type of ion present. Water partially splits up into ions {this is why it's such a good solvent for ionic compounds}. It splits into hydrogen ions and hydroxide ions. H2O → H+ +OH- So at the cathode there will be two ions present: the metal ion and the hydrogen ion from the water. Which element is actually produced at the cathode depends on how reactive the metal is. If the metal is very reactive, such as potassium or sodium, then it is unlikely to be discharged. Hence hydrogen will be produced. If the metal is unreactive such as silver, the metal will be produced. To work out which ion "wins", the metal or the hydrogen, compare their reactivities in a the reactivity series. The one that is most reactive, will not be produced at the cathode. A similar situation occurs at the anode. Hydroxide ions {from the water} are usually discharged at the anode ultimately producing oxygen. However, if the concentration of the ions of Halites (group 7) are much higher than that of the hydroxide ions, then the halite ions are discharged. Sulphates are never discharged. OH- → OH + e- 4OH → 2H2O + O2. Q9) Sodium chloride is dissolved in water and subjected to electrolysis. Explain what you see at each of the electrodes. Answers | «Advanced static electricity | Circuits» 

नमस्ते! Hindi is an Indo-European language spoken as a first language in majority states in northern India and second language in many countries where these people have emigrated. It is written with the Devanagari, script which fairly closely follows the phonetics of the language. Spoken Hindi is very similar to spoken Urdu — as such they are both often classified as part of the Hindustani language. Feel free to use Wiktionary's Hindi Category to study words. Basic Hindi in 1,2,3 Steps. "To be reorganized" 

Axiom of choice: If formula_1 is a surjective map, then there exists a map formula_2 such that formula_3 is the identity (trivial) map. Lemma: Every set can be well-ordered. 

An Application Service Provider (ASP) is a business model for delivery of business or IT services across a network. It may also be a service enterprise whose primary business is delivery of application services using the ASP model. ASP usually implies a central (vendor) data center that runs the application. For further discussion, see the Wikipedia article. 

Solving linear inequalities involves finding solutions to expressions where the quantities are "not" equal. A number on the number line is always greater than any number on its left and smaller than any number on its right. The symbol "&lt;" is used to represent "is less than", and "&gt;" to represent "is greater than". For example:  -5 -4 -3 -2 -1 0 1 2 3 4 5 From the number line, we can easily tell that 3 is greater than -2, because 3 is on the right side of -2 (or -2 is on the left of 3). We write it as formula_1 (or as formula_2). We can also derive that any positive number is always greater than negative number. Consider any two numbers, "a" and "b". One and only one of the following statements can be true: Now we can go on to solve any linear inequalities. Solving Inequalities. Solving inequalities is almost the same as solving linear equations. Let's consider an example: formula_29. All we have to do is to subtract 4 on both sides. We will then get formula_30, and that is the answer! Note, however, what you get is not a single answer, but a "set" of solutions, i.e., any number that satisfies the condition formula_30 (any number that is less than 9) can be a solution to the inequality. It is very convenient to represent the solution using the number line:  &lt;-------------------o  6 7 8 9 10 11 Let us try another more complicated question: formula_32. First, you may want to expand the right hand side: formula_33. Then we can simply rearrange the terms so that all the unknown variables are on one side of the equation, usually the left hand side: formula_34. Hence we can easily get the answer: formula_35. This solution is represented on the number line below. Note that the solution requires a closed circle ("●"), because the formula_11 is greater than "or equal to" 4.  &lt;-+-----+-----+-----+-----+-----+--&gt;  -6 -5 -4 -3 -2 -1 Inequalities with a variable in the denominator. For example consider the inequality In this case one cannot multiply the right hand side by formula_38 because the value of x is unknown. Since x may be either positive or negative, you can't know whether to leave the inequality sign as formula_39 (ie less than), or reverse it to &gt; (ie greater than). The method for solving this kind of inequality involves four steps: Compound Inequalities. A compound inequality is a pair of inequalities related by the words "and" or "or". In an "and" inequality, both inequalities must be satisfied. All possible solution values will be located between two defined numbers, and if this is impossible, the compound inequality simply has no solutions. Consider this example: formula_56 and formula_57. First, solve the first inequality for x to get formula_35. All "and" inequalities can be rewritten as one inequality, like this: formula_59formula_57 (write x between two ≤'s or &lt;'s or both with the smaller number on the left and the larger number on the right). Now, we can graph this inequality on a number line as a line segment. Remember, all solutions to ≤ or ≥ must be graphed with closed circles. Interpret this graphic as "all numbers between -4 and 2, including -4 and 2."  &lt;-+-----+-----+-----+-----+-----+--&gt;  -6 -4 -2 0 2 4 Now, let us consider "or" inequalities. "Or" inequalities usually do not have a set of solutions that satisfies both. Instead, they usually have two sets of infinite numbers that are solutions to each one. Because of this, "or" graphs define which numbers satisfy either equation. For example: formula_61 or formula_62. First, solve for x in the second inequality to get formula_63. Now, graph the two inequalities on the same number line. Remember to use open and closed circles accordingly.  &lt;-------------o ●--------&gt;  -1 0 1 2 3 4 Solving Inequalities with Absolute Value. Since formula_64 A inequality involving absolute value will have to solved in two parts. Solving formula_65 The first part would be formula_66 which gives formula_67. The second part would be formula_68 which solved yields formula_69. So the answer to formula_65 is formula_71  &lt;-+-----+-----+-----+-----+-----+-----+-----+-----+-----+-----+-----+-----+--&gt;  0 1 2 3 4 5 6 7 8 9 10 11 12 Graphing Linear Inequalities. The graphing of linear inequalities is very similar to the graphing of linear functions. A linear inequality is written in Previous: Solving equations &lt;br&gt; Next: Quadratic functions 

JLPT Level N5. The following is a list of 81 kanji, which is most of the kanji necessary to pass the N5 level of the Japanese Language Proficiency Test (JLPT), prior to the updating of the list several years ago. There are now 103 kanji in the level N5 exam (24 more than the kanji listed below). Kanji used in the N5 test are used in more difficult levels, too. As of 2010, the JLPT no longer publishes an official vocabulary list. Unofficial lists for all JLPT levels are however published on . 

Active Server Pages (ASP) is a technology that creates dynamic Web Pages by executing Java Script or VB Script statements on a server, generating "HTML", and sending the resulting Webpage to a browser. ASP is packaged as a part of the "Personal Web Server (PWS)" and "Internet Information Server (IIS)" suites. 

About the platform. macOS is the primary operating system for the Macintosh computer. It was originally a system designed privately by Apple Inc, however with Mac OS X, it has been based on Unix. Specifically, a modified FreeBSD operating system called "Darwin". There are many different kinds of software that can be developed for Mac OS X. People generally think of applications, but we'll briefly cover some of the other kinds. Types of Software for Mac OS X. Applications. Applications are what people generally think of when they think about software for Mac OS X. Cocoa applications include: Finder, Mail, Address Book, Safari, Microsoft Word, and Microsoft Excel. Anybody can develop applications using Apple's free development tools which includes XCode. Mac OS X applications are developed using Objective-C though there are other possible programming languages that could be used. The most popular languages for use on the macOS platform is Objective-C which could be thought of as Mac OS X's "native language" since the Mac OS X libraries, or "frameworks", all have an Objective-C interface. Objective-C includes everything that plain C can do, and adds object-oriented programming. See: . C++ can be used in developing for the Mac, but generally, it is used in addition to Objective-C rather than being in place of Objective-C. Using both Objective-C and C++ is called "Objective-C++" and is considered to be optional when developing software for Mac OS X: See for a lesson on the basics of Objective-C may also be of assistance. Some preliminary thoughts: Objective-C is the language most commonly used in Mac OS Programming. Objective-C entered Mac OS X and has ancestry in NeXT. . Before you learn Mac programming you "must " know the basics of C since it is the basis for Objective-C. There used to be three separate APIs for developing a Mac application with a GUI: 1. Classic (Mac OS 9 and lower). Developing for the Classic API is no longer done. When Mac OS X first came out, users and developers had a huge investment in software written for Mac Classic OS and Mac OS X used to have an emulation mode so that users could run their old software. Apple has long since stopped support of the Classic API and Classic emulation in Mac OS X. 2. Carbon (Mac OS 8.5 up to and including Mac OS X 10.6 Snow Leopard). Carbon was an API for developers to update their applications that used the Classic API to be run without the Classic emulator. Carbon was a great way that Apple provided developers to upgrade their software to run on Mac OS X without having to totally rewrite their software, but Carbon, like Classic, is no longer supported by Apple. 3. Cocoa (All versions of Mac OS X). Cocoa is the most native API that can be used to develop applications for Mac OS X that are truly "Mac-like". Generally, Objective-C will be used along with Cocoa, though there are other options such as Cocoa-AppleScript and Cocoa-Python, but Cocoa-Objective-C is really the "mainstream" way to develop Cocoa applications. Resource Forks Files in Mac OS X have a feature that is unique to Mac OS and that is that each file on disk can have two "forks". This feature used to be used for Classic and Carbon applications to separate code from resources (such as menus, windows, etc.), and the Mac OS X file system still supports two forks, but you should only use the "data fork". The resource fork is non-standard and can be lost when transferring Mac files to other file systems. AppleScripts. Another "native language" for developing Mac OS X applications is AppleScript. AppleScript is a language that Apple invented to automate repetitive tasks. The AppleScript application is located on your Mac at /Applications/Utilities/AppleScript Editor. AppleScript can be used to record AppleEvents, the events that applications send to themselves or to other applications. Why don't you try it out. Open AppleScript Editor, press the record button, do some things with your other applications and watch the script write itself. AppleScript can be used alone or it can be used along with XCode to develop Cocoa Applications using mostly AppleScript instead of Objective-C. This option is mostly for experienced AppleScript programmers who don't know Objective-C. Automator Workflows. Apple also provides an application called "Automator" that can be used to easily automate repetitive tasks. It is located at /Applications/Automator.app Shell Scripts. Mac OS X has an application called Terminal that provides a command-line interface to Mac OS X. It is possible to develop scripts for the command line. Terminal.app is located at /Applications/Utilities/Terminal.app To create a shell script, you need a text editor. There is a text editor that comes with Mac OS X called "TextEdit.app". It is located in /Applications/TextEdit.app. But actually, what is better than TextEdit is a program such as TextWrangler.app which is available for free from the following link: http://www.barebones.com/products/textwrangler/ The shell that Terminal.app uses by default is called "bash". Here is a simple tutorial on developing bash scripts http://www.maclife.com/article/columns/terminal_101_automate_terminal_bash_scripts We won't go any more deeply into shell scripts here in this wikibook, but it's just good to know what they are. You can always google for more information now that you know what to google for. Command Line Tools. When you open Terminal and you learn how to type in commands. The commands are usually command-line tools or scripts. Above, we just talked about developing your own scripts with a text editor. It's also possible to develop your own command-line tools, using XCode. This is an advanced thing to do. Usually, power-users will write a shell-script (or some other kind of thing such as an AppleScript or an Automator Workflow) but it's good to know what a command-line tool is. Command-line tools have a textual user-interface rather than a graphical user interface (GUI). Java. Java used to be treated by Apple as a "first class language" to develop for Mac OS, however in recent years, Apple has less support for Java. Now with Mac OS X 10.7 "Lion" and 10.8 "Mountain Lion", Java doesn't even come pre-installed in Mac OS X. Java is still available, but users have to download Java from Oracle's website and install it themselves. Apple's Mac App Store doesn't even allow Java apps to be sold at their store calling Java "deprecated". However, there still are Mac developers who use Java because it has the advantage of being cross-platform compatible. For example, the same source-code can be used to generate software that runs on Mac, Windows, and Linux. Apple has said that Java reduces the Mac to the "least common denominator". That's why they support it less. Python. Python is somewhat supported by Apple. In fact, Python is shipped with Mac OS X and is part of the System Folder. There are third-party libraries that allow developers to develop applications using Python and Cocoa together, but these are not very well maintained, and Python on the Mac is most suitable for developing command-line utilities, or cross-platform scripts that aren't really very Mac-like. Ruby. Similar to Python. Websites. Most Mac users use Safari for their web browser. Safari uses the standards set by w3c.org You can develop websites that work with Safari by following the standards of the w3c.org. Remember to validate your HMTL, CSS, and JavaScript. HTML Validator: http://validator.w3.org/ CSS Validator: http://jigsaw.w3.org/css-validator/ JavaScript Lint: http://www.javascriptlint.com/online_lint.php If you're developing websites using your Mac and using Safari, remember to test your webpages on other platforms and with other web browsers. Mac OS X Specific Languages. Objective-C is really the "native" language for Mac OS X development You could call AppleScript a "native" language too, but it isn't really used to make commercial applications. It was designed to be used by real power-users to automate their tasks. Although it is possible to use AppleScript to build Cocoa applications in XCode, this would be more for users who already know AppleScript and don't want to learn Objective-C. 

This part contains a survey of the world of computer hardware. This chapter introduces computer hardware and provides a very short history of computing. This section introduces the computing environment thrugh a brief history of computers and an overview of the subjects within hardware. Outline of computer history. There is some disagreement about when the digital age began. After all, the Babylonians used a kind of abacus about 500 BC. If we call it the computer age, we can get a much cleaner mark. 1951 was a watershed year, and two events identify it as the beginning. The UNIVAC became the first commercially available digital computer. That same year a research team at MIT completed the first user-interactive computer (named Whirlwind) that used a keyboard and a television screen or CRT. Considering developments by decades gives a pretty clear picture of the evolution of computers. Vacuum tubes. computer hardware notes Computer Hardware. All notes of Computer hardware 

Sanskrit is an old Indo-Aryan language once spoken throughout India. Many poems, epics, and prayers are written in it. Although not widely spoken anymore, it gave rise to many of the modern languages of India and influenced several more languages that are not related to it; Sanskrit is also still used as a language of culture and religion in India. In the past, Sanskrit was written in a variety of scripts by people from various places. Today, a modified version of the Devanagari script has been agreed upon as a standard. Sanskrit script. Harvard-Kyoto coding. a A i I u U e ai o au aM aH (or a:) ka kha ga gha Ga (nasal "ng") ca cha ja jha Ja (nasal "ny") ta tha da dha na Ta Tha Da Dha Na (retroflex) pa pha (not "fa") ba bha ma ya ra la va ha See also. Wikner's introduction to Sanskrit is a good start. "Sanskrit: An Easy Introduction to an Enchanting Language" by Professor Ashok Aklujkar is also recommended. 

This is an elementary Latin course accompanied with a detailed grammar based upon Kennedy's Public School Latin Grammar designed to introduce one to the world of . A basic understanding of would be helpful; however, it is not required. Basic definitions of terms will be explained in Lessons 1 and 2, and later elaborated as needed. For detailed explanations and examples of English grammatical terms, please consult the English Grammar textbook. However, Latin grammar is quite different from that of English, and thus it requires different grammatical terms to explain the concepts. These will be taught as needed. Preface. This book will attempt to teach the reader about Latin from the ground up. Please read the Introduction to the origins and structure of Latin carefully, as it introduces the concept of a stem. As is typical in many other languages, the infinitive stem (present tense, active voice) is used for conjugating verbs. [The introduction of additional information in parentheses is done simply to avoid confusing a student who has already had exposure to Latin.] Advices. Parts of this book may have been edited by people who do not speak English as their first language. All Wikibooks are written in the particular English dialect of the writer, which may not be standard usage. If you see something particularly unclear, please feel free to correct it, but please alter this article in a constructive manner. If something doesn't make sense to you: if you are skilled enough then delete or better fix it else try to emphasize that text marking it with some keyword, e.g., "grade 12 [grade] American [system] revert?" reporting that thing in the of the book. The "revert" keyword allows your editors to know that you are not a skilled editor but that you are just trying to learn, and you are confused. Your changes are not permanent. In any case, please, edit this book responsibly. Spoken Latin. This is a test chapter to teach those who wish to learn Latin which they can use in their daily lives. About the Book. Please leave ideas for additional chapters on the page. See also. Cultural insights in other wiki books / wiki articles. 

=What is Latin?= Latin was the language originally spoken in the region around the city of Rome called Latium. It gained great importance as the formal language of the Roman Empire. All — including Italian, French, Spanish, Portuguese, Romanian, and others — descend from a Latin parent, and many words in English and other languages today are based on Latin roots. Moreover, Latin was a "lingua franca", the learned language for scientific and political affairs in Europe, for more than one and a half thousand years, being eventually replaced by French in the 18th century and English by the middle of the 20th. Latin remains the formal language of the Roman Catholic Church to this day and is the official language of the Vatican. Romance languages are not derived from Classical Latin, a literary language for writing and oration, but rather from Vulgar Latin, the language spoken by the common people, or " vulgus, " of Rome. Classical Latin and Vulgar Latin (Romance) differ (for example) in that Romance had distinctive stress whereas Classical had distinctive length of vowels. In Italian and Sardo logudorese, there is distinctive length of consonants and stress, in Spanish only distinctive stress, and in French even stress is no longer distinctive. Stress refers to the emphasis of pronunciation on syllabic units. Most English nouns not derived from other parts of speech have an emphasis on the first syllable. Foreign loan words in English sometimes retain their original stress, which may be on the second or third syllable, though assimilation into English will usually result in a vowel shift towards emphasis on the first syllable. Another major distinction between Classical and Romance is that modern Romance languages, excluding Romanian, have lost their case endings (suffixes at the end of the word used in place of prepositions) in most words (some pronouns being exceptions). Romanian is still equipped with several cases though some, notably the ablative, are no longer represented. Here are some current English words which are Latin derivativesː =Introduction to the Latin Language= Simple and Compound Words. In Latin, words are either: Word Parts. Inflected words (i.e., words having ending- or spelling-changes according to their grammatical functions in the sentence) have a stem and a root. The Stem The stem is the part of the word to which various suffixes are added. The final suffix determines either the role of the word in the sentence (for example, when a Roman slave wished to address his "dominus" (master), he used the vocative form "domine" -- equivalent to "O master" in English) or the person/subject involved in the action (for example, "I dominate" may be expressed as "domin-or", and "they dominate" as "domin-antur"). In these cases, "domin-" is the stem and "-us", "-e", "-or" and "-antur" are suffixes. The addition of such suffixes is called "inflection". This is discussed further in the Summary. The Root The root is the part of the word that carries the essential meaning. For example the stem of "agitō" (I drive onward) is "agit-", whose root is "ag" (do, drive), which is in common to words of similar meaning: "agō" (I do, drive), "agmen" (that which is driven, such as a flock), etc. Notice the essential difference between a root and a stem. To the root "ag" has been added a suffix "(i)tō-" which denotes frequency of action (so "agit-" means to do or drive more than once, hence "agit-ō", I agitate, I keep (something) moving, I urge, I impel). In contrast, English uses word order more than inflection to determine the function of a word within a sentence. English also uses words like pronouns (I, she, etc.) and prepositions (to, at, etc.) where Latin generally prefers inflexions. Thus "dom-ī" (noun -- "at home"), "ag-unt" (verb -- "they do/drive"). Primitives Primitives occur when both the stem and the root are the same. For example, in the word "agere" (to do, drive) both the stem and the root are the same: "ag-". Derivatives Derivatives occur when the root or stem is modified. For example, the stem "flamm-" from the noun "flamma" has the root "flag" ("blaze"), "nōscō" (I know) from the verb "nōscere" has the root "gnō-" ("know"). Suffixes Latin attaches suffixes ("endings") to stems to turn them into words (most stems and roots cannot be used in sentences without an ending). This inflection is essential to forming Latin sentences. The various suffixes and their translations will be learned in the later lessons. =Types of Words used in Latin= Nouns. A noun (Latin: "nōmen") is "something perceived or conceived by the mind." There are two kinds of nouns: Substantives and Pronouns. 1. Substantive ("nōmen substantīvum") is a name simply denoting something perceived or conceived: "psittacus" - the parrot, "nix" - the snow, "virtus" - virtue. 2. Pronoun ("prōnōmen") is a word used in place of a "substantivum", usually when the "substantivum" is already known: "ea" - she, "ille" - that man Nouns have changing endings on the stem (known as declension) and three incidents: number, gender and case. Number concerns whether the thing referred to is singular or plural (and the ending shows this); gender classifies a substantive as masculine, feminine or neuter (this determines how the endings of adjectives and pronouns behave) and case (where the ending must show how the noun fits in to the sentence). Adjectives and Pronouns must agree in all incidents when they refer to a substantive. Verbs. Verbs ("verba") express an action or a state of being, e.g., "agō" (I do), "dīxit" (he said), "venīs" (you come). Conjugation is the term for adding inflections to verb stems to indicate person (first, second or third), number (singular or plural), tense (present, future, imperfect, perfect, pluperfect or future perfect), voice (active or passive), and mood (indicative, subjunctive or imperative). A verb can be either "finite" or "infinite": 1. Finite verbs ("verba finīta") are inflected and have a subject, e.g., I run, you run, he runs, they drive, the computer is turned on. 2. The infinite verbs ("verba infinīta") are not inflected and have no subject, e.g. to run, to drive, to turn on, to have drawn. "Participles", which are inflected as substantives rather than as verbs, may also be considered infinite, e.g., the "running" boy. Modifiers. 1. Adjectives ("adiectīva") are used to describe nouns. They indicate a quality perceived or conceived as inherent in, or attributed to, something denoted. E.g., "vir magnus" (the great man), "puella pulchra" (the fair girl) 2. Adverbs ("adverbia") are similar to adjectives, except that they are used to qualify verbs, adjectives or other adverbs, rather than nouns. In practice, they restrict the meaning of the verb or adjective by specifying how or how much. E.g., "currō celeriter" (I run quickly), "pugnat fortiter" (he fights bravely), "vērē iūcundus est" (he's really nice), "incrēdibile callida est" (she's incredibly clever). Other. Particles are uninflected words that provide extra meaning. 1. Prepositions ("praepositiōnēs") are little words which tell you how one word is behaving in relation to another word ("the duck was near the pond", "she went towards the wood"). In Latin, the noun that follows a preposition takes a particular ending (called a "case"), depending on the nature of the relationship, or on the nature of the preposition itself. E.g., "ad" (by), "in" (in), "sub" (under). What all this means is that a preposition is a sort of adverb, telling you how something is done. For example, "you go" is a simple statement, but "you go in" suggests that you don't just "go", you go so as to enter something, and so you need a noun for the "something". In English, we might say "you go into the house". In Latin, this would be: "in domum inīs". Notice the form "in domum", which means "into" the house -- you're going into it, you're not yet exactly inside it (the ending -um of "domum" is called "accusative"). When you are inside the house, what you do is "in" the house, which is "in domō" (the ending -ō of "domō" is called "ablative"). 2. Conjunctions ("coniunctiōnēs") join together clauses and sentences. E.g., "et" (and), "atque" (as well as), "sed" (but). 3. Interjections ("interiectiōnēs") are exclamations used to express feeling or to gain attention. E.g., "ō!" (oh!), "ēheu!" (alas!), "ecce!" (behold!). Articles. Latin has no definite article or indefinite article, respectively "the" and "a/an". When translating Latin into English the appropriate article must be added. =Summary= =Exercises= 

Spelling and Pronunciation. The Latin alphabet, on which the English alphabet is based, has mostly the same letters as the English alphabet, except that it has no &lt;k&gt; or &lt;w&gt;, and that in its original form, it lacked &lt;j&gt;, which only some modern texts use, and &lt;u&gt;. Many European languages use the Latin alphabet as the basis for their own alphabet. Latin pronunciation has varied somewhat over the course of its long history, and there are some differences between Old Latin, spoken in the Roman Republic, Classical Latin, spoken in the Roman Empire, Medieval Latin, spoken in the Middle Ages, and Ecclesiastical Latin, spoken in the Catholic Church. This text focuses on the pronunciation of Classical Latin. Note that Latin, as written by the Romans, did not have &lt;j&gt;, &lt;k&gt;, &lt;u&gt;, &lt;w&gt;, or macrons over vowels (the lines indicating that vowels are long), although they did sometimes mark long vowels with apices (e.g. &lt;ó&gt; for /oː/); macrons are used today as pronunciation guides and do not necessarily need to be written. /w/, /ʊ/, and /uː/ were all represented with &lt;v&gt;. Modern texts often use &lt;v&gt; for /w/, &lt;u&gt; for /ʊ/, and &lt;ū&gt; for /uː/. In some modern texts (this Wikibook not included), &lt;j&gt; is used for /j/. Declension Tables. The following tables will be both referenced and explained in all of the following sections, and hence are placed here. Note that nouns in the 3rd declension nominative can have any ending, hence why none is given in bold. Grammar Part 1: Nouns and Their Role in Sentences. Nouns in Latin are inflected, which means that endings (also known as suffixes or "suffices") are appended to the end of the stem to denote these things: Most nouns in English can be modified to indicate number (cat versus cats), and many pronouns can be modified to indicate case (who versus whose) or gender (he versus she, his versus hers). Case is especially important in Latin as meaning cannot be determined by word order as it can be in English, but purely by word endings, or "inflection". Indeed, the words in a Latin sentence can appear in almost any order with little change in meaning. Two sentences with the word orders "Sam ate the orange" and "The orange ate Sam" could potentially mean the same thing in Latin, though the spellings of "orange" and "Sam" would have to change slightly to denote which was the subject (the one eating) and which was the object (the one being eaten). It is important to note here that although the genders of many words make sense (for example, "puella", meaning a girl, is feminine) many are simply assigned and hold no real meaning. Luckily, as you will find, the gender can often be determined by the spelling of the word (words ending in "us" are almost always masculine, and words ending in "a" are almost always feminine). For many words, however, you will simply have to memorize their gender. Adjectives themselves must match the number, case, and gender of the noun (be it a substantive or a pronoun) they modify. If a noun is nominative singular feminine (see case table below), then the adjective describing it must also be nominative singular feminine. If the noun is accusative plural masculine, then the adjective must be accusative plural masculine. This will be expanded on in the Adjectives section below. The advantage of this system is that adjectives do not need to be adjacent to their respective nouns, as one would be able to tell which noun they modify by which noun they appear to agree with. Declension. All substantives are part of one of 5 categories, called declensions. A substantive is a stem, modified by adding a declension suffix. Each declension has a set of standard suffixes that indicate case and number. Usually gender is indicated by the suffix, although there are many exceptions. Therefore, you must memorize the gender of every substantive you learn. By familiarizing yourself with the above tables, you could deduce that originally the suffix indicating number, case, and gender was the same for every noun. However, as the language developed, nouns with a common stem formed declensions and sounds changed. Similar processes happen continually over time, even today. The above tables allow you to familiarize yourself with the existence of each declension, though by no means are you expected to memorize it now. Nonetheless, you will have to memorize it as you are formally introduced to individual cases and declensions in future lessons. Because of its introductory purpose, it is considerably simplified and incomplete, and therefore should not be used as a reference in the future. Adjectives are also classed into declensions which must match the declension of the noun they describe: Pronouns are not part of any declension, as they are all irregular, and simply have to be memorized. Case. Cases (Latin: "casus") determine the role of the noun in the sentence in relation to other parts of the sentence. There are six cases, Nominative, Genitive, Dative, Accusative, Ablative, and Vocative. Vocative case (Lesson 3), can be considered a sort of miniature case, generally not being accepted as a true one. Additionally, some nouns have a vocative case, which will be covered later. As nominative and accusative are the most basic, these will be taught first (the rest will be covered in later lessons). Gender. All substantives, including inanimate objects, have a particular gender (genera), which is either masculine, feminine, or neuter. For example, Vir, "a man," is masculine. Marītus, "a husband," is also masculine. Puella, "a girl," is feminine. Māter, "a mother," is feminine. Even inanimate objects are assigned gender, including all the moons, stars, trees, tools, and so forth. Logic will give you little help in determining what the genders of inanimate objects are, and with many nouns memorization is required. Luckily, for many nouns, the spelling of the word indicates the gender. Certain rules may be utilized to determine the gender of an inanimate substantive. Declension is a good indication of gender, especially for 1st and 2nd declension substantives. 1st declension substantives (substantives with an -a suffix) are usually feminine and second declension nouns (substantives with an -us suffix) are usually masculine or neuter. There are a few exceptions, and they will have to be learned. 3rd declension nouns can be either masculine, feminine or neuter (thus the gender will often have to be memorized). 4th declension nouns are usually masculine, sometimes neuter while 5th declension nouns are usually feminine. Nouns undeclined, words which are not substantives but used as such, sentences used as substantives and the products of trees are generally neuter. 1st/2nd declension adjectives alternate the set of endings depending on the gender of the noun it describes (see the next section below). If the adjective describes a feminine noun, the adjective must use 1st declension endings, if the adjective describes a masculine noun, the adjective must use 2nd declension masculine endings, if the adjective describes a neuter noun, the adjective must use 2nd declension neuter endings. 3rd declension adjectives use the same set of endings for masculine and feminine nouns. However, a slightly different set of endings are used when describing neuter nouns. Adjectives. As stated above, adjectives must match the gender, number, and case of the noun (be the noun a substantive, or a pronoun) they modify. Similar to the "Sam ate the orange" example above, if the adjective uses the wrong declension it could change the meaning of the sentence. For example, "The girl loves big trees," versus "the big girl loves trees" have different meanings. There are many occasions where logic cannot be used to determine the gender of inanimate objects, as genders are generally arbitrary when the noun has no literal gender. Furthermore, the declension of the noun, often determined by the spelling, can in turn be used to determine the gender, especially for the 1st and 2nd. However, this is never the case for the third declension, as the declension itself is not primarily assigned to any gender and the spelling of the nominative ("default") stem is random, leaving you with no hints. A noun and its adjective must also be in the same case. Otherwise, it is impossible to tell which nouns pair up to their respective adjectives in a sentence, as the words in a Latin sentence can appear in any order. See the examples below. Recapitulation. Therefore: Before you proceed to the next lesson, complete the exercises below so you will be able to apply this knowledge to Latin. 

The Nominative Case. The Nominative case refers to the subject of a sentence. For example: The girl is pretty "The girl" is the subject of this sentence. In its simplest form a sentence will have a subject stated as a noun and will give some further information about the subject. The second part of this sentence tells the reader that the girl is pretty. This is called predicating the noun. This sentence consists of a subject and a predicate. As you know from English, an adjective is a word that denotes some quality, which in this sentence is attractiveness. The noun and adjective are joined together by the word "is", which is called the copula. Note that the copula simply connects the words and gives almost no information about the subject. The sentence in Latin has the same grammatical elementsː puella est pulchra The noun is followed by the predicate. The only difference is the absence of an article which has to be supplied by the translator. Puella can be translated as "girl", "the girl", or "a girl". Can you tell which word is the copula? Translate the followingː Which region of Europe was the Roman historian Tacitus referring to as Caledonia in his book "Agricola", which records the military campaigns of his father-in-law? Translate the followingː Note the conjunction given in the Vocabulary, and translate the followingː Give the meaning of the complete word on this inscription fragment from Roman Britainː Vocabulary. Key to Vocabulary: Overview of Adjectives. An adjective is any word that qualifies a noun. For example: Adjectives in Latin. Adjectives must agree with the nouns they describe in gender, number, and case. These words will look like the adjective antiquus (old, ancient): Third declension adjectives typically look more like ferox, ferocis (wild, bold). This is because the third declension has no stem assigned to the nominative singular. Adjectives often come after the word they describe. Since word order is not central to the meaning of a Latin sentence the adjective may appear anywhere in the sentence. In the following examples the "-us" is masculine (m.), "-a" is feminine (f.) and "-um" is neuter (n.). So magnus is masculine, magna is feminine, and magnum is neuter. Basic verbs. Verbs in Latin work quite differently than those in English. Study the following table then view the examples below. Personal Endings. Archaic Latin was spoken and written in Europe for over two thousand years and since all languages change gradually this sometimes makes it difficult for beginners to see patterns of change. English has also had a long development that is now divided into three periods called Old English, Middle English and Modern English. Compare the following English verbsː The contraction of the archaic "laborao" to "laboro" would have undergone the same gradual process. The archaic "amao" (I love) eventually became "amo". If you look at the Vocabulary you will see that "amat" and "amant" retain the original letter "a" of the stem. Further Examples. Example 2. Notes: In the same way, the adjective "pulcher -ra -rum" must agree with "puella" in gender, number, and case, so the correct form is "pulchra" (agreement with the feminine nominative singular noun of the first declension). Example 4. Notes: The adjective "magnus -a -um" in this case must agree with "lūdī" in gender, number, and case, so the correct form is "magnī" (masculine nominative plural). Third Declension Nouns and Adjectives. Third declension nouns and adjectives follow a different pattern. The nominative singular stem is not defined, and as such, any letter (or letters) can serve as a third declension stem. For example, "Māter" (mother) is a third declension noun in the nominative case. When pluralized, it becomes "Mātrēs". "-ēs" is attached to the end of a third declension noun to pluralize it, as opposed to changing the ending completely, because there is no uniform way to do so. You may have also noticed that that the "e" in "Māter" was dropped when pluralized. This often happens when a stem is attached to a third declension noun of similar spelling (example, "Pater" (father) becomes "Patrēs") Examples: Third declension nouns are listed with the nominative case and the genitive case to provide the main stem. For example: All other types of nouns are also generally listed with the genitive Adjectives with a nominative ending in -is and the same stem in the nominative and in the other cases (eg. fortis) end in -e in the neuter and -ia in the neuter plural. For example: 

English abbreviations derived from Latin. Common English abbreviations from Latin. Latin External Links. ^ Latin ^ 

Authors of the Latin Wikibook include the following contributors: and many anonymous Wikibooks contributors. 

Latin This page provides a list of Latin phrases and their English translations. This page is copied from the article . Check that page to see the latest changes to this page. V. from , the Free encyclopedia. ^ Latin ^ 

Grammar: The Accusative. As you learned in the last lesson, the verb 'esse' (to be) usually takes the nominative case, because then the word after it is a complement. Most other verbs take the 'accusative' case. In a sentence, the accusative is the "what" - in English grammar, this is known as the direct object. For example: The girl sells the box. What did the girl sell? The box. Thus, box is the direct object, and when we translate it into Latin: Cistam, then, is in the accusative, because it is the direct object. Again, when an adjective describes a noun in the accusative case, the adjective must agree in number, case, and gender. Because Latin uses cases to mark the subject and the object of a sentence, word order does not matter. Consider: Examples of Adjectives Agreeing with the Nominative and Accusative Case. "Bonus", a first and second declension adjective, is masculine, nominative, and singular to agree with "puer", the word it is describing. "Ferocem", a third declension adjective, is masculine, accusative, and singular to agree with "canem". "Canem" is accusative because it is the object of "amat". Here is an example of plural adjectives: The words "bonus" and "ferocem" become "boni" and "feroces" to agree with the plurals "pueri" and "canes". However, if a girl (puella) happened to love that boy: "Bonus" must become "bona" in order to modify "puella", which is feminine. Finally, if the girl isn't good, but rather wild: Even though "puella" is first declension, "ferox" remains third declension. In the same way, a good lion would be "bonus leo". Exercise 2. Determine whether the adjective agrees with the substantive in all three categories: case, gender, number. Grammar: The Use of the Accusative. The newly introduced verbs, ama-t, curri-t, and porta-t take the accusative as the 'object'. Unless specified, any verb you look up in the dictionary will take the accusative, not the nominative. This means that they are transitive verbs, verbs that happen to someone or something, e.g.:  I heal you. ("acc.")  You make my day. ("acc.")  She hit your arm. ("acc.") In the examples above, the bold words are the subject of the sentence clause. Because something happens "to" them, they can't be in nominative. 

A noun, or noun substantive, is a part of speech (a word or phrase) which functions as the head of a noun phrase. The word "noun" derives from the Latin "nomen" meaning "name", and a traditional definition of nouns is that they are only those expressions that refer to a person, place, thing, event, substance, quality, or idea. They serve as the subject or object of a verb and as the governed term of a preposition, and can co-occur with articles and attributive adjectives. There are different groups of nouns: Each of these different groups of nouns have different properties, each making them different in how we use them. Thus, nouns are names of objects, places, people and things. They are used with adjectives to describe something, and with verbs to show an action. Concrete nouns. Concrete nouns are proper nouns and common nouns. Proper nouns. Proper nouns are the names of people, places, groups or dates: as, "Adam", "Boston", "the Hudson", "the Romans", "the Azores", "the Alps". They almost always have a capital letter as their first letter. Example: Common nouns. Common nouns are the names of a sort, kind, or class, of beings or things: as, "beast", "bird", "fish", "insect", "creatures", "persons", "children". They often refer to objects or things which we can see, touch and feel, like the word "chair". Example: Individual nouns. Their refer to only one thing of the same kind, for eg: man, player, cow, chicken, minister. Collective nouns. Collective nouns are the names of a groups of objects or many individuals together: as, "council", "meeting", "committee", "flock". Example: Abstract nouns. Abstract nouns are the names of some particular qualities considered apart from its substance: as, "goodness", "hardness", "pride", "frailty". They are often names of the things that we cannot touch or see, but are there all the same. Example: Verbal nouns. Verbal nouns or participial nouns are the names of some actions, or states of being; and are formed from a verb, like a participle, but employed as a noun: as, Sui generis. A thing sui generis, (i.e., of its own peculiar kind,) is something which is distinguished, not as an individual of a species, but as a sort by itself, without plurality in either the noun or the sort of thing: as, "galvanism", "music", "geometry". Inflections of Nouns. Nouns have modifications of genders, numbers, and cases. Genders. Genders, in grammar, are modifications that distinguish objects in regard to sex. There are three genders; the masculine, the feminine, and the neuter: Hence, names of males are masculine; names of females, feminine; and names of things inanimate, literally, neuter. Numbers. Numbers, in grammar, are modifications that distinguish unity and plurality. There are two numbers; the singular and the plural. The singular number is that which denotes but one: as, The plural number is that which denotes more than one: as, Regular plurals. The plural form is usually represented orthographically by adding "s" to the singular form. The phonetic form of the plural morpheme is by default. When the preceding sound is a voiceless consonant, it is pronounced . Examples: "boy" makes "boys"; "girl", "girls"; "chair", "chairs"; "cat", "cats". Where a noun ends in a sibilant sound, the plural is formed by adding "es" (pronounced ), which is spelled "es" if the word does not already end with "e": "glass" makes "glasses"; "dish, dishes; witch, witches; phase, phases; judge, judges". Most nouns ending in "o" preceded by a consonant also form their plurals by adding "es" (pronounced ): "hero" makes "heroes"; "potato, potatoes; volcano, volcanoes". Nouns ending in a "y" preceded by a consonant drop the "y" and add "ies" (pronounced ): "cherry" makes "cherries"; "lady, ladies". Proper nouns (particularly those for people or places) ending in a "y" preceded by a consonant form their plurals regularly: "Harry" makes "Harrys"; "Germany, Germanys". This does not apply to words that are merely capitalised common nouns: as, "P&amp;O Ferries". A few common nouns ending in a "y" preceded by a consonant form their plurals regularly: "henry" makes "henrys"; "zloty, zlotys". Words ending in "ey" form their plurals regularly, in order to avoid the unpleasant-appearing vowel sequence "eie": "monkey, monkeys". Almost-regular plurals. Many nouns of Italian or Spanish origin are exceptions to the oes rule: "canto" makes "cantos"; "piano, pianos; portico, porticos; quarto, quartos; solo, solos". Many nouns ending in a voiceless fricative mutate that sound to a voiced fricative before adding the plural ending. In the case of changing to the mutation is indicated in the orthography as well: "calf" makes "calves"; "bath, baths; mouth, mouths; house, houses". Some retain the voiceless consonant: "proof" makes "proofs"; "moth, moths; place, places; dwarf, dwarfs or dwarves; hoof, hoofs or hooves; staff, staffs or staves; turf, turfs or turves; roof, roofs or rooves". Irregular plurals. There are many other less regular ways of forming plurals. While they may seem quirky, they usually stem from older forms of English or from foreign borrowings. Irregular Germanic plurals. The plural of a few Germanic nouns can also be formed from the singular by adding "n" or "en", stemming from the obsolete weak declension: "ox" makes "oxen"; "child, children". The plural is sometimes formed by simply changing the vowel sound of the singular, in a process called umlaut (these are sometimes called "mutated plurals"): "foot" makes "feet"; "goose, geese; louse, lice; man, men; mouse, mice; tooth, teeth; woman, women". Some nouns have singular and plural alike, although they are sometimes seen as regular plurals: as, "aircraft, sheep, deer, fish, cod, trout, head, cannon". Generally, plurals refer to several species or kinds of animal, while the unmarked plural is used to describe multiple individual animals; one would say "the classification of fishes", but "five fish in an aquarium". Irregular plurals of foreign origin. Such nouns often retain their original plurals. In some cases both forms are still vying: for a librarian, the plural of "appendix" is "appendices"; for physicians, the plural of "appendix" is "appendixes". A radio engineer works with "antennas" and an entomologist deals with "antennae". The "correct" form is the one that sounds better in context. Correctly formed Latin plurals are the most acceptable, in academic and scientific contexts. In common usage, plurals with "s" are sometimes preferred. Cases. Cases, in grammar, are modifications that distinguish the relations of nouns or pronouns to other words. There are three cases; the nominative, the possessive, and the objective. The nominative case. The nominative case is that form or state of a noun or pronoun, which usually denotes the subject of a finite verb: as, The subject of a finite verb is that which answers to who or what before it: as, Boy is therefore here a noun in the nominative case, or nominative. For example: The possessive case. The possessive case is that form or state of a noun or pronoun, which usually denotes the relation of property: as, Boy is here a noun in the possessive case, or possessive. The possessive case of nouns is formed, in the singular number, by adding to the nominative "s" preceded by an apostrophe; and, in the plural, when the nominative ends in "s", by adding an apostrophe only: as, singular, boy's; plural, boys'; sounded alike, but written differently. The objective case. The objective case is that form or state of a noun or pronoun which usually tells the object of a verb, participle, or preposition: as, The object of a verb, participle, or preposition, is that which answers to whom or what after it: as, Boy is therefore here a noun in the objective case, or objective. The nominative and the objective of nouns, are always alike in form, being distinguishable from each other only by their place in a sentence, or by their simple dependence according to the sense. For example: The declension of nouns. The declension of a noun is a regular arrangement of its numbers and cases. Thus: A short syntax. The subject must be in the nominative case, as "You say it." The subject is placed before the attribute, as "Peace dawned on his mind," except the following cases: a question, as "How many loaves have you?" imperative mood, as "Go you," strong feeling, as "May she be happy!" a supposition, as "Were it true," "neither" or "nor", as "Neither shall you touch it," emphasis, as "Here am I," no regimen, as "Echo the mountains round," dialogue, as "My name is Hassan," and the adverb "there", as "There lived a man." A noun in apposition is put in the same case as the noun it explains, as "But he, our gracious master, knows us." A possessive is governed by the name of the thing possessed, as "Man's life." A possessive comes immediately before the governing noun, as "Nature's peace," except the following cases: an intervening adjective, as "Flora's earliest smells," affirmation or denial, as "The book is not John's," a possessive without sign, as "Brother Absalom's house," or "David and Jonathan's friendship." The predicate is governed by attribute in objective case, as "I found her." A noun or a pronoun put after a non-transitive verb or participle, agrees in case with a preceding noun or pronoun referring to the same thing, as "The child was named John." The case of absolute noun or pronoun depends on no other word, as "Your fathers, where are they?" 



Elemental analysis is the process for determining the partial or complete chemical formula for a substance. Most commonly, it involves the complete combustion in air or oxygen of the substance and then quantifying the amount of elemental oxides produced. In the case of organic compounds, the carbon is converted to carbon dioxide and the hydrogen to water. From these, the percent carbon and percent hydrogen in the substance can be found and compared with a proposed chemical formula for the substance at hand. Element test: Put a small amount of the solid into a small piece of Na metal then roll it around the solid, followed by introduction into a fusion tube. The tube is heated with a gentle flame at a slow rate (in order not to evaporate N2 present in solid) then strong heating till the bottom of the tube become red hot. The tube is then put in a beaker containing a minimal amount of water then heated, cooled, filtered and the filtrate divided into three parts. 1. Test for nitrogen: The filtrate and ferrous sulfate are boiled and cooled and dilute sulfuric acid is added. If green or blue color occurs the solid contain nitrogen. The chemistry behind what happened:&lt;br&gt; Na+C+N --&gt; NaCN&lt;br&gt; FeSO4+NaCN gives Fe[CN]2&lt;br&gt; Fe[CN]2+4NaCN give the complex Na4[Fe(CN)6]&lt;br&gt; ferrous oxidizes to ferric by the acid so 3Na4[Fe(CN)6]+4Fe3+ --&gt; Fe4[Fe(CN)6]3 2. Test for sulfur: The filtrate is exposed to dilute acetic acid and lead acetate, yielding a brown or black precipitate. The chemistry behind what happened:&lt;br&gt; Na+S --&gt; Na2S&lt;br&gt; Na2S +Pb(CH3CO2H)2 yields lead sulfide, a black precipitate. or: The filtrate and sodium nitro prusside yield a violet color Na2S+Na2[Fe(CN)5NO] --&gt; Na4[Fe(CN)5NOS] = violet color 3. Test for chlorine: a) In the absence of N or S: The filtrate is exposed to dilute nitric acid and silver nitrate. Formation of a white precipitate suggests the presence of chlorine. b) In the presence of N and/or S : The filtrate is exposed to dilute sulfuric acid then boiled to 1/3 initial volume and cooled. Formation of a white precipitate after the addition of dilute nitric acid and silver nitrate suggests the presence of chlorine. Equations: NaCN+AgNO3 --&gt; AgCN white ppt&lt;br&gt; Na2S+AgNO3 gives Ag2S black ppt&lt;br&gt; In the presence of N or S these two precipitates may interfere with the white colour of the result of the chlorine test. Therefore dilute sulfuric acid is added because in the presence of N or S: Na2S+ dilute H2SO4 gives H2S gas NaCN+ dilulte H2SO4 gives HCN gas There is no interference with the white colour expected from the chlorine test in solid. 

Chromatography involves the physical separation of a mixture of compounds. Chromatography can be used as a purification method but also sees wide use for the identification of compounds based on their chromatographic behavior. =Theory= There are many variations of chromatography, but all involve the dissolution of an analyte into a fluid known as the mobile phase and the passage of this fluid solution across a stationary phase, often a solid or liquid-coated solid. As the mobile phase comes into contact with the stationary phase, some of the analyte molecules dissolve or adsorb onto the mobile phase. The more the molecules of that substance are retained, the slower their progress through the chromatographic apparatus. Different substances will then move through at different rates, ideally resulting in distinctly identifiable retention times for each substance. Commonly used chromatographic techniques are identified through the nature of the stationary and mobile phases used, the method for passing the mobile phase through the apparatus, and how separated components are identified. =Paper chromatography= In paper chromatography the stationary phase is a specialized paper made to absorb water to a high level. The mobile phase is usually water or a concentrated salt solution. Paper chromatography has many uses in forensic chemistry due to it's simplicity and availability. However, paper chromatography is limited by the characteristics that only water soluble components can be separated and inaccuracy in RF values. This makes paper chromatography mostly useful to distinguish the differences between two residues rather than their similarities. =Thin layer chromatography= In thin layer chromatography (TLC) a plastic or glass plate is coated with the stationary phase, often alumina, silica, or alkylated silica. The analyte is dissolved in a quick-drying solvent and spotted near the bottom of the plate. The edge of the plate beneath the spot or spots is then dipped and left in a solution of the mobile phase, either an organic solvent or aqueous solution (depending on the nature of the analyte and stationary phase). Capillary action is then allowed to draw the solvent front through the spotted analyte, carrying with it and in the process separating out the analyte's components. =Gas chromatography= In gas chromatography (GC) the analyte and mobile phase must both be gases or be readily introduced into the gas phase by heating. The mobile phase gas must be inert and not reacting with the sample to be analysed. Examples of inert gases are helium and nitrogen gas; while not as inert as helium and nitrogen, hydrogen may also be used. The gases are passed through a long, narrow (and most often, coiled) tube either packed with a porous stationary phase or whose inner walls are coated with a stationary phase, and the analyte components are detected as they emerge from the far end of the tube. The tube is commonly known as GC column. Often a time-varying temperature gradient, from lower temperature to higher temperature, is applied to the tube. This first allows the analyte components to partition into the stationary phase and then, as the temperature rises, to differentially force them back into the mobile phase. Common detectors for gas chromatography are flame ionization detector (FID), electron capture detector (ECD) and mass spectrometry (MS). Different types of sample analysis would require the use of a different type of detectors. =Column chromatography= Column chromatography, like gas chromatography, uses a tube packed with a stationary phase, but the mobile phase is a liquid instead of a gas (It is sometimes known as liquid chromatography or LC). Instead of temperature gradients, a gradient in the composition of the liquid phase can be used to separate components. Column chromatography can be performed on larger molecules which may not be readily introduced into the gas phase. On the other hand, because of the increased viscosity of liquids compared with gases, liquid chromatography can be a more ponderous process. HPLC (variously "high-pressure liquid chromatography" or "high-performance liquid chromatography") speeds the process and improves its selectivity and sensitivity to a significant degree by forcing the mobile phase through the chromatographic column with high-pressure pumps. =Detection methods= The root of the word chromatography, "chroma" (Greek khrōma, color) and grafein is "to write", indicates that the separated components in some forms of the technique can be identified by their color alone. But chromatography has now long been performed on colorless compounds that can be identified in other ways. Analyte components on thin-layer chromatography plates are often identified under ultraviolet light, or by chemical staining in, for example, an iodine chamber or potassium permanganate. Gas chromatographic analytes are detected by changes in the ionization levels of a flame at the output end of the column or by changes in the electrical conductivity of the gas mixture at the end of the column. Liquid chromatography fractions are often analyzed through spectrophotometric techniques, notably UV-visible spectroscopy. When separation with GC or LC is performed in tandem with mass spectrometry (the "hyphenated" techniques of GC-MS and LC-MS), masses of individual fractions are rapidly determined. These methods are frequently employed in analytical and forensic science. 

Chapter 1 Botany as a Science. Botany is the branch of biology concerned with the scientific study of plants. Traditionally, botanists studied all organisms that were not generally regarded as animal. However, advances in our knowledge about the myriad forms of life, especially microbes (viruses and bacteria), have led to spinning off from Botany the specialized field called Microbiology. Still, the microbes are usually covered in introductory Botany courses, although their status as neither animal nor plant is firmly established. Plants are living entities, and material presented within "Biology" will have relevance here, most particularly at the cellular and subcellular levels of organization (Chapter 2). Both plants and animals deal with the same problems of maintaining life on planet Earth — their approaches seem quite different, but the end result is the same: continued existence in an organized state, as part of a universe whose tendency is towards greater disorganization. Back on Earth, however, it is a fact that microbes, plants, and animals comprise a very interdependent system. We divide them apart, because our minds work best that way. We categorize and learn common features or properties of the categories. This approach is neither right nor wrong, but is clearly efficient for our minds. Nonetheless, it is desirable to regularly step back and realize that the boundaries between categories are often just constructs, and exceptions to our categories usually abound. It was alluded to in the opening definition that Botany is a science. Just what makes Botany, or anything else a science? It is important to acquire a grasp of the fundamentals of science itself to fully appreciate both how botanical knowledge was gained as well as how it can be used. It usually becomes uninteresting to acquire facts simply for the sake of knowing. Humans do not just appreciate mountains "because they are there", they climb them because they are there! Living Systems. Biology is defined as the study of life, and Botany is that discipline within Biology concerned with the study of living organisms called plants and with certain other living things that have been traditionally studied with plants i.e. Fungi and Algae. Defining 'Plant'. Like many words in common usage that apply to biological entities or concepts, the term plant is more difficult to define than might be at first obvious. Although botanists describe a Kingdom Plantae, the boundaries defining members of Plantae are more inclusive than our common concept of a "plant". We are tempted to regard "plant" as meaning a multicellular, eukaryotic organism that generally does not have sensory organs or voluntary motion and has, when complete, a root, stem, and leaves. However, botanically only vascular plants have a root, stem, and leaves, and even some vascular plants, such as certain carnivorous plants and duckweed, fall afoul of that definition. But to be fair, the vascular plants are the plants we tend to encounter every day and that most people would readily regard as "plants". A more significant point of departure between Plantae and plants occurs among the seaweeds. Technically, only a relatively minor group of seaweeds (the chlorophytes or green algae) are members of the Kingdom Plantae. The majority of seaweeds, like the kelps (very large brown algae from the Order Laminariales), despite a superficial appearance of such, lack true stems, leaves, roots, and any kind of vascular systems as found in higher plants. Thus, the kelps are not Plantae; but are they plants? Certainly if we regard the green algae as plants, it is difficult to exclude the more prominent red and brown algae of our coastal waters. Another, much broader definition for "plant" is that it refers to any organism that is photoautotrophic—produces its own food from raw inorganic materials and sunlight. This is not an unreasonable definition, and is one that focuses on the role plants typically play in an ecosystem. However, there are photoautotrophs among the Prokaryotes, specifically photoautotrophic bacteria and cyanophytes. The latter are sometimes called (for good reasons) blue-green algae. Then there arises the problem that many people would consider that a mushroom is a plant; a mushroom is the fruiting body of a fungus (Kingdom Fungi) and not "photoautotrophic" at all, but "saprophytic". However, there are more than a few species of flowering plants, fungi, and bacteria that are not autotrophic, but "parasitic". We cannot hope to offer a firm answer. The list of characteristics that separate the Plantae from the other biological kingdoms provides at least a technical definition, but realize it is only a technical definition. The problem this lack of precision or agreement in the definition of "plant" presents is one of understanding statements, often encountered in "Wikipedia" (and other) articles, of the sort: "...xylem is one of the two transport tissues of plants". In general it cannot be assumed this means all plants, algae through flowering plants. It very probably does not include fungi or bacteria. Indeed, it is usually safest to assume the discussion is about vascular plants (essentially the ferns, conifers, flowering plants, and a few others; see discussion below on "General Terminology") unless stated differently (e.g., "...in vascular and non-vascular plants this is "such and such). Plants as Organisms. The distinction between life and non-life is not as easily made as you might think. There exist intracellular "parasites" that are progressively less alive in terms of being metabolically active. Plants and their Uses. There can be no disputing the fundamental significance of plants to the ecology of our planet. Photosynthetic plants utilize energy arriving from the sun to create complex organic molecules from inorganic substances, and by this process contribute oxygen to the atmosphere. Advanced animal life is very much dependent upon this source of oxygen, as well as the organic molecules that form the basis of nearly every food web on the planet. However, humans utilize plants in many ways, especially as sources of pleasure, food, and material for shelter, clothing, and more. Consider here the role plants play in our everyday lives and in our economy. Introduction to Plant Classification. At the beginning of this chapter it was suggested that each of us categorizes information we encounter on a daily basis. Our minds seem to want to find relationships between facts and observations, to erect mental bins in which to place new items with previous "facts". This natural human process is the basis for prejudice, in as much as "facts" categorized together can become strongly associated. But these are personal constructs. In order for scientists of many races, speaking many languages, and coming from all manner of backgrounds and experiences to work productively together to solve common problems, the objects with which they work must be classified within a universally accepted framework. The classification of living things is called systematics, or taxonomy, and ideally should reflect the evolutionary history (phylogeny) of the different organisms. Taxonomy arranges organisms in groups called taxa, while systematics seeks clues to their relationships. The dominant system of "Scientific Classification" is called Linnaean taxonomy, and includes classification ranks as well as an organism naming convention called binomial nomenclature. Traditionally, all living things were divided into five kingdoms: However, this five-kingdom system has been replaced by Carl Woese's three-domain system, which focuses on phylogenic roots and comparison of DNA structures. The older approach utilized visual observation as the basis of classification. The three domains reflect whether cells have nuclei (eukaryotic) or not (prokaryotic), as well as differences in cell membranes and cell walls. Recall (and review as necessary) how these groupings relate to the sequence of events in the evolutionary history of life as summarized in . You will return to the subject of Scientific Classification to consider in much more detail the groups of organisms studied in Botany, beginning with Chapter 7. First, however, we shall turn our attention to the structure and function of cells and eventually to gain an understanding of plant structure ("plant anatomy") and function ("plant physiology"). General Terminology. In Section II of this text we will delve much deeper into "plant" systematics. But you should be aware of some general terms related to classificatory schemes that are used regularly in discussing plants. You have probably encountered these terms many times, although may not be aware of their exact definitions. For example, much of the material in Section I of this textbook is biased towards flowering plants. That is, much of the descriptive material here as well as at Wikipedia refers specifically to these. Flowering plants are angiosperms; plants that have flowers and produce seeds, and comprise the majority of the plants we would normally encounter in say a nursery if not on the street, field, or empty lot. Seed-bearing plants include both the angiosperms and the gymnosperms, the latter now treated as a modern group called conifers. The conifers are also common plants, especially in higher latitudes, but bear cones instead of flowers. Both conifers and flowering plants develop vascular tissues internally that conduct fluids (especially water) throughout the plant. Included in the vascular plants are ferns. Ferns have vascular tissue, but reproduce by spores. They do not produce seeds and do not bear flowers. 

Chapter 2 Chapter 2. Plant Cells Introduction. A cell is a very basic structure of all living systems, consisting of protoplasm within a containing cell membrane. Only entities such as viruses— on the boundary between non-living chemicals and living systems—lack cells or basic cell structure. All plants, including very simple plants called "algae", and all animals are made up of cells, and these are organized in various ways to create structure and function in an organism. Biologists recognize two basic types of cells: prokaryotic and eukaryotic. "Prokaryotic cells" are structurally more simple. They are found only in single-celled and some simple, multicellular organisms (all bacteria and some algae, which all belong to Bacteria and Archaea domains). "Eukaryotic cells" are found in most algae, all higher plants, fungi, and animals (Eukarya domain). Thus, differences between these two cell types are critical to how an organism is classified, and an important consideration in the evolutionary sequence of life on the planet Earth. Plant Cell Structure. Nearly all cells are too small to be seen with the unaided eye. As always there are some exceptions, but generally magnification is required to detect a cellular structure. In plants, a good hand-lens or loupe (see photo at right) will sometimes suffice, but in working with cells or observing how cells are organized to form tissues and structures, a high power is used. Questions: Basic Cell Function. You should, by now, have a general appreciation for the complexity of cellular structure. Improvements in microscopy, especially development of the , have revealed that cells are not merely membranous sacks containing fluid of gel-like consistency. The degree of organization of the cytoplasm into organelles and their membranes should have you convinced that much (perhaps most) of what is really going on around you on this planet is occurring at a scale that is simply inaccessible to your eyes. And while you cannot be expected to directly observe chemical reactions at a molecular scale, contemplate that you cannot, even with powerful optics, directly observe most of the structure where these reactions are somehow controlled to produce outcomes favorable to life—indeed, are life. Hopefully, as you acquire knowledge and become a biologist—a botanist—you will learn to recognize the relevant phenomena by their macroscopic expressions (that which you can readily observe with the unaided eye). To appreciate basic cell function, it is necessary to first list the processes or outcomes that cells must accomplish to further existence. More specialized functions will be discussed under plant cell structure, as our interest must eventually focus on plants. For now, recall that in your reading you have already encountered these several basic abilities of cells: Now explore each in turn. Think initially of a single-celled organism with no special abilities, only a "will" to stay alive and perpetuate itself. Remember, the environment will not be kind. The cell must grow and reproduce to counter the tendency of outside forces to breakdown molecular structure and destroy life. Then consider the situation where a cell is part of a multicellular organism, and may be performing more limited and specialized functions. Questions: Plant Cell Specializations. We will learn about the cells of algae and other organisms (e.g., bacteria and fungi) traditionally covered within Botany in later chapters on those organisms (Chapters 5 - 7). Here, we concentrate on the cells of plants. The simplest type of plant cell is called a parenchyma cell and most of the basic metabolic and reproductive processes of the plant occur in these cells. A term for "parenchyma" cells with chloroplasts, is chlorenchyma cells. Other plant cell types that we shall be considering are: Laboratory Exercises for Chapter 2 »&lt;br&gt; Discussion of questions for Chapter 2 »&gt; Energy. How does plant cells get its energy? It gets its energy nutrients. With those nutrients, they gather the energy through sunlight. They use photosynthesis to convert materials into energy for the plants to power its cells. 

Chapter 4 Chapter 4. Plant Vegetative Organs &lt;br&gt; Introduction. As was noted in the previous chapter, most plant cells are specialized to a greater or lesser degree, and arranged together in tissues. A tissue can be "simple" or "complex" depending upon whether it is composed of one or more than one type of cell. Tissues are further arranged or combined into organs that carry out life functions of the organism. Plant organs include the leaf, stem, root, and reproductive structures. The first three are sometimes called the "vegetative organs" and are the subject of exploration in this chapter. Reproductive organs will be covered in Chapter 5. The relationships of the organs within a plant body to each other remains an unsettled subject within plant morphology. The fundamental question is whether these are truly different structures, or just modifications of one basic structure (Eames, 1936; Esau, 1965). The plant body is an integrated, functional unit, so the division of a plant into organs is largely conceptual, providing a convenient way of approaching plant form and function. A boundary between stem and leaf is particularly difficult to make, so botanists sometimes use the word shoot to refer to the stem and its appendages (Esau, 1965). The Leaf. The plant leaf is an organ whose shape promotes efficient gathering of light for photosynthesis. The form of the leaf must also be balanced against the fact that most of the loss of water a plant might suffer is going to occur at its leaves (transpiration). Leaves are extremely variable in terms of their size, shape, and adornments (such as small hairs on the face of the leaf). Although the leaves of most plants carry out the same basic functions, there is nonetheless an amazing variety of leaf sizes, shapes, margin types, forms of attachment, ornamentation (hairs), and color. Examine the Leaves (forms) page to learn the extensive terminology used to describe this variation. Consider that there are functional reasons for the modifications from a "basic" type. The Stem. The stem arises during development of the embryo as part of the "hypocotyl-root axis", at the upper end of which are one or more cotyledons and the shoot primordium. The Root. The root is the (typically) underground part of the plant axis specialized for both anchoring the plant and absorbing water and minerals. Basically, there are two types of roots normally spotted for plants grown on ground namely : taproot and fibrous root Most of the material you have read discusses the root organ as found in the angiosperms (flowering plants). However, among the vascular plants, only Psilotales lack such an organ, having instead rhizomes that bear hair-like absorbing structures called rhizoids (Eames, 1936 in Esau, 1965). Questions:  4-1. At this point the conceptual differences between cell types, tissues, organs,  and organisms may be somewhat confusing. Using the leaf as an example, describe  this structure in a way that considers the cell types, tissues, and organs for  that part of the leaf where photosynthesis is concentrated. 

GCSE Science/Electricity You've probably done quite a lot of work in basic electrical circuits in lower school, but we will revise them here so don't worry if you've forgotten them. Before we start you need to know a whole host of circuit symbols. Q1)See how many of the following you can name: You should already know about half of them. For the others see here. You need to learn all these symbols by heart, but don't worry about that just yet. It's much easier to learn the symbols if you know what the components actually do, and we will be looking at that later on. What is a circuit? A circuit is a loop in which electricity can flow. Consider a simple circuit with a cell, a switch and a bulb. The current flows from the cell, via the closed switch to the bulb. The bulb lights because the electricity carries energy. But even after the bulb there current still needs to go around the rest of the loop back to the cell. If there is a break "anywhere" in the loop, the bulb will not light. This is a concept that is not all that difficult, yet most people do not understand it. Ask a selection of adults you know and many will say:&lt;br&gt; "Electricity flows from the cell to the bulb, where it gets used up" Q2) What is wrong with the above statement? Q3) Look at the selection of circuits below. What is wrong with each of them? What we have been looking at so far has been largely revision. You should already be familiar with what a circuit is and how it works. Let's now go on to some more advanced work. « Electrolysis | Current, voltage resistance and Ohm's law 

Contributors. While Wikibooks offers somewhat clearer opportunities for "authorship" than Wikipedia, there remains the fact that anything put here is really just a contribution, and everyone who furthers the effort is a contributor. In this respect there are no "authors". If you are interested in developing a particular subject within the field of botany as part of this textbook, please start the module and link it here. We will work out the chapter arrangement after you get started. Of course, you can also just expand on an existing chapter. «Return to Contents Page 

Did you get many right? You should already know, cell, battery, bulb, switch and resistor. You will have to learn the others as well before the exam. «Back 

Chapter 5 Chapter 5. Plant Reproduction Vegetative Reproduction. Vegetative reproduction is a type of asexual reproduction—other terms that apply are "vegetative propagation," "clonal growth", or "vegetative multiplication". Vegetative growth is enlargement of the individual plant; vegetative reproduction is any process that results in new plant "individuals" without production of seeds (see "The Seed" below) or spores. It is both a natural process in many, many species as well as one utilized or encouraged by horticulturists and farmers to obtain quantities of economically valuable plants. In this respect, it is a form of cloning that has been carried out by humankind for thousands of years and by "plants" for hundreds of millions of years. Sexual Reproduction. The Flower. The flower is the reproductive organ of plants classified as "angiosperms"—that is, the flowering plants comprising the Division Magnoliophyta. All plants have the means and corresponding structures for reproducing sexually, and these other cases will be explored in later chapters. However, because flowering plants are the most conspicuous plants in almost all terrestrial environments, we justifiably devote this chapter to the flowering plants alone. You will learn how other plant groups (and non-plant groups, such as fungi) reproduce sexually in Section II of "The Guide". The basic function of a flower is to produce seeds through sexual reproduction. Seeds are the next generation, and serve as the primary method in most plants by which individuals of the species are dispersed across the landscape. Actual dispersal is, in most species, a function of the fruit: structural parts that typically surround the seed. But the seed contains the germ of life of the next generation. The Seed and Germination. the primary purpose of the seed is one of preserving the continuity of life—starting a new generation in a new physical location. For large plants (shrubs and trees), this can be especially important because successful germination and growth close to the parent may be difficult or impossible; the established plant monopolizes light and water resources in its immediate vicinity. Seeds can also serve the function of overwintering or surviving harsh conditions. The entire generation—every individual—may die in the Fall or the dry season. In many annual species, only the seed exists during unfavorable dry or cold conditions. The Fruit. The fruit is the actual agent of dispersal in most flowering plants. 

Electronics GCSE Science/Electricity What is electricity anyway? So far in this module we have been using words like electricity, current, voltage, and resistance without actually explaining in depth what these words mean. On this page you will learn what these words mean, and how they relate to one another. What is electricity on an atomic level. Let's think about a wire, made of copper. You should already know that it is made of particles called copper atoms, which have a positive nucleus surrounded by negative electrons. We need not concern ourself with the nucleus, because it does not move. They are stuck in a rigid crystal structure. Instead we need to look at the electrons. Copper, like all metals, has some loose electrons. These electrons are not held rigidly by the atoms but are free to roam about anywhere as long as they stay somewhere in the metal. We call them "conduction electrons" because they are the ones that conduct electricity. To see how electric current flows imagine a small petri dish, with an even smaller one glued inside. The channel between the two dishes is then filled with peas. The channel represents the wires in a circuit, the peas the electrons. Charge is the quantity of free charged particles (in this case electrons). The SI unit of charge is the coulomb. You "could" count the electrons, but there are an awful lot of electrons! Each electron carries a "tiny" charge, of −1.602×10-19 coulomb. We can turn this figure on its head and say that if we grouped together 6.2415×1018 electrons we would have −1 coulomb of charge. Electric field. An electric field is a region where an object experiences a force due to its charge. The strength of the electric field (electric field strength E) is described by considering how much force is experienced by a unit charge (a charge of 1 coulomb) when it's placed in the field. The electric field at any point is therefore expressed in newtons per coulomb. Some fields (such as the field between two parallel charged plates) have the same field strength throughout the field - uniform fields. Others (particularly radial fields due to isolated point charges) do not have constant field strength. In a battery there is a negative and positive terminal. The conventional current flows from positive to negative (the large line to the small line). Because opposites attract, these charges in the battery will be attracting each other, but they can not move directly to each other through the battery because of the chemical processes. If there is a complete external circuit, this attraction from the battery will give the free electrons in the metal (i.e. the wire) a force which will make them move. If you think about it these electrons are being forced to the other side of the battery because of this attraction. We call this driving force the Electro-motive Force (e.m.f). Emfs are measured in volts, and are sometimes referred to simply as voltage. The larger the emf (the voltage) the more quickly electrons flow round the circuit. What is the rate of flow of the electrons?...the current. so, a larger voltage means a larger current. Current. Now imagine that you were to put your finger on one pea and push it in a clockwise direction. All the peas would move because they are all touching one another. This is just what happens when an electric current flows. Current is the flow of charged particles (the particles are usually electrons). The SI unit of current is the ampere (symbol A). To understand how current is defined think of standing at a given point in the wire. Electrons are flowing past you. One Ampere is a flow of one coulomb going past every second. Definition of the ampere. BUT this flow of charge idea is NOT the definition of the ampere. The ampere is defined in terms of the force produced between two wires each carrying identical currents:- "The ampere is that constant current which, if maintained in two straight parallel conductors of infinite length, of negligible circular cross-section, and placed 1 meter apart in vacuum, would produce between these conductors a force equal to 2 x 10-7 newton per meter of length." Voltage. Having looked at charge and current, we now need to look at what voltage means. As you know, electrons repel each other. If you hold one electron near another electron, you have to push against it to hold it in place. If you try to bring it even closer, you have to do work (force times distance) to get it to that new position. The voltage between one point and another is simply how much work per coulomb is required to move any small test charge from point A to point B. In most electronic components, it doesn't matter much which path the test charge takes in-between point A and point B. Voltage is related to the "energy" of the charges. Let's go back to the peas in the petri dish. They can be pushed slowly or they can be pushed quickly. The faster they go, the more energy they have. It's a similar situation for the electrons, except the push isn't provided by a finger! It's provided by the battery. The battery gives the charges energy. This energy is given up to the various components in the circuit, e.g. bulbs, resistors etc. The energy per unit charge is called the voltage (or the potential difference). Definition of the volt. One volt means one joule of energy per coulomb of charge. More accurately it has 2 definitions: Electromotive Force is the amount of energy converted from non-electrical to electrical form when driving 1 coulomb of charge around a completed and closed circuit. Potential Difference is the amount of energy converted from electrical to non-electrical form when driving 1 coulomb of charge around a completed and closed circuit. The potential difference between 2 points in a conductor is defined as 1 volt, if 1 Joule of energy is converted from electrical form to non-electrical form, when 1 coulomb of charge per second (1 amp) flows through it. Note: This will only occur between 2 points in a conductor, that has a resistance, defined as 1 ohm. See resistance below. Resistance and Ohm's law. Resistance is easy to understand. It's just how difficult it is for the charges to flow through an electrical component or from one point to another in an electrical circuit. Imagine a group of walkers travelling down a road. They approach a fork in the road. To the left is a flat straight road leading to a nearby town. To the right is a huge mountain, over which a steep and winding road traverses. This road also leads to the nearby town. Naturally all the walkers chose the left route. Let's now suppose that there are millions of walkers. They are jam packed on the road, and they are all in a hurry to get to the town as quickly as possible. Now when they come to the fork in the road which way should they go? Most will still go to the left, but a few might chose to go to the right, the road is more difficult but there is no traffic jam, so they might get there quicker. It's a similar thing with moving charges. Like charges repel, so they would rather not pack in very closely together. Some routes, like wires have very low resistance, while other routes like bulbs have much higher resistance. More charges will go down a low resistance route than a high resistance one. The unit of resistance is the ohm, represented by the Greek letter Omega (Ω). Ohm's Law. This law relates resistance, current, and voltage. It's very easy to remember because it's obvious when you think about it. Let's think of a wire carrying current from a battery, to a bulb then back to the battery. The voltage of the battery provides the energy of the flowing electrons. Let's assume we want to increase the rate of flow of charge. Remember that current is the rate of flow of charge so we want to increase the current. So formula_1 Rearranging gives : In symbols: This is Ohm's law. Learn it by heart. Q2) A wire has a resistance of two Ω and a voltage of 3 V is applied. What is the current? Q3) A wire has a 5 V potential difference applied. the current is 10 A, what is the resistance? Q4) A current of 2 A is flowing through a wire with a voltage of 20 V. What is the resistance of the wire? (Remember to give the units for all your answers). Answers | «Circuits | Circuits Part 2» 

« Contents Page Introduction to the Botany Study Guide. This "Study Guide to the Science of Botany" is a textbook at Wikibooks shelved at the and intended to establish a course of study in the subject of Botany, utilizing articles provided in Wikipedia, with links to other relevant web sites and other Wikibooks as appropriate. In some cases, portions of the text from "Wikipedia" articles have been used to materially develop introductory text within the Guide. For the new user, it need be pointed out that "Wikipedia" differs from a standard encyclopedia in two important respects: 1) it is a hypertext document, and 2) it is open and editable, and therefore constantly changing. For the student following this or any guide through "Wikipedia" to cover a specific subject, it is recommended that each article (page) be read first in its entirety, before any hyperlinks are followed to other topics or explanations. It is too easy, otherwise, to simply become lost in a maze of links, and miss the main thrust of an article presented as an assignment from the Guide. Because "Wikipedia" is constantly changing (and, it is believed, improving) the quality of each article encountered will be variable. Some articles are well written and go to adequate depth, whereas others, lacking a proponent, are shallow and incomplete. Short or sloppy looking articles may contain questionable facts. These short-comings should diminish with time, but can be a problem for the student. One clear advantage to using this Guide linked to a hypertext like "Wikipedia" is the "circular redundancy with serendipity" factor that arises when an article is read and its hyperlinks followed; this factor can be a powerful learning tool. The persistent reader is subjected to a fairly high degree of repetitive reading, often presenting slightly differing perspectives on the same general topic, with the result that learning comes from redundancy and seeing difficult concepts presented in more than one way. At the same time, some hyperlinks lead down less relevant paths, bringing new and unanticipated knowledge. If, as a student, you are truly interested in mastering the subject of botany, you must be prepared to read beyond the basic assignments; in some cases, beyond Wikipedia to explore other, "outside" web sites. It seems likely that the typical user of the "Study Guide to the Science of Botany" is not necessarily an active student taking a course in botany at the high school (AP) or college level, but a person with a strong interest in plants—an amateur naturalist or a gardener. Therefore the guide must incorporate both the basic biological and physiological aspects of plants as well as extensive taxonomy-based coverage of the diversity of plants and related organisms. The amount of material now available on the web covering the latter subject is becoming nothing short of phenomenal. In effect, one now has access to much of the world's plant diversity, with photographs and descriptions, in many cases from web sites maintained by specialists. One goal of the guide is to provide a systematics-based approach to capturing this kind of information, hopefully giving the student a strong background in plant systematics. The importance of this approach is not that everyone should become a taxonomist—or become more familiar with plant taxonomy, a specialized field of botanical science with a relatively narrow following—but that appreciation for (and understanding of) species diversity is most critical at this time in our earth's history marked by accelerated species extinctions and destruction of native ecosystems by both human population expansion and man-induced spread of non-native species. The "Study Guide to the Science of Botany" includes two other "parallel" documents intended to enhance the usefulness of the Guide. These could also be used separately or independently as source documents for a beginning course in Botany. They are the "Discussion" pages and the "Laboratory Exercise" pages. Both are explained in detail in the next Section titled "How to use this Guide". 

« Contents Page How to Use the "Botany Study Guide" The purpose of the "Study Guide to the Science of Botany" is to weave - out of the information on Life Science and especially Botany contained in "Wikipedia" - a course of study for the student or layman. It is anticipated that this course will be either supplemental to instruction being received at a school or college, or will be self-directed. In either case, the Guide is not a novel and should not be approached as one. A smooth flow of dialogue is simply not possible and should not be anticipated. The Guide may be closer to the sometimes disjointed notes generated by a student from a lecture or careful reading of a detailed textbook. Within each subsection of a Chapter, introductory text is followed by one or more "reading assignments" of the form: Following (that is "clicking" on) the link (to "Wikipedia" "Botany" in this case) will open an article intended to provide the details of the Chapter subsection. Recommended articles should be read from top to bottom, and then re-read following some or all of the links embedded in the article to other articles for expanded elucidation or to clarify terms; that is, in most cases, completion of an "assignment" (recommended article) includes at least some or all articles linked to the first. Obviously, it cannot be the case that all links are followed to articles, whose links are then followed to articles, and so on until no new material is encountered. It is likely there would be no quick end to such a pursuit. The amount of time spent wandering beyond the original article is partly a personal matter of how much the reader is getting out of the foray than anything else. Realize it is certainly possible to wander well off the subject at hand. As in the example above, notes are provided with assignments giving some direction for pursuing links. An instruction NOT to follow links simply means the additional material will be encountered later in the course of instruction, and going beyond the assigned article may provide too much detail for a beginning student. The following example: specifies that two other links ARE part of the assignment. Other links encountered may be followed to expand your knowledge or, as always, to aid in understanding of technical terms encountered. Hyperlinks included with the text in the Guide are there simply for convenience, usually to topics somewhat peripheral to the main one. In all cases, finding your way back to the Guide may become tricky, but we have to leave this up to you to establish, beyond pointing out that your browser's "Back" button is intended for this purpose. Discussion Questions. At the end of each subsection are posted one to several questions. In general, you will get more out of these questions if you write out your answer on a piece of paper. You may wish to accomplish this on the re-read, allowing each question to guide your quest for an answer. A discussion page for each chapter provides answers to the questions posed. However, the questions are intended to be thought-provoking, and may not have a single straight-forward answer. Answers on the discussion pages are also necessarily much longer than would be expected of any one student; it is expected that each student answer will fit somewhere within the broad discussion presented. Laboratory Exercises. A natural sciences course laboratory unit is supposed to provide hands-on experience in exploring topics raised in the text and lecture units. The best that a website can give towards this goal is a manual that is liberally provided with pictures and diagrams. The student must provide the "hands on" from the neighboring natural world. Fortunately, in botany, this is much easier to accomplish than in almost any other field of science. Both the outdoors, the local market, and (if available) a botanical garden can be sources of materials for study. Indeed, we may teach the structure of a "pome" using an apple in the hope that the student will end up with a pear. In using any of the Laboratory Exercises, it is always best to read through the entire module before actually doing anything. Resist the temptation to view the material as an instruction manual to be followed in a specific order. For one thing it is difficult to write a module that covers, at each step, all that the student should know before proceeding on to the next step. The value of any exercise will be significantly enhanced if you have a pretty good idea where it is going in advance. General Navigation. The "Study Guide" is divided into Sections and Chapters which define the subject material of each module. At the bottom of each text page (main text of a module), is a short version of the Table of Contents, allowing the reader to jump between chapters within a Section. Here is an example of the "Wiki Contents Table" for Section I:  Note that at the beginning of each module, links are provided to both the previous chapter and the succeeding chapter, as well as to the main Table of Contents. Links to units associated with a module, as for example to a Laboratory Exercise, appear near the end of the module Final Note. As a final note, read the next Section and consider how you might make a contribution to the Guide. 

Further Discussion. The questions posed in Chapter 3. Plant Structure are discussed further on this page. Remember, some questions are intended to be thought-provoking and more than one answer may be "correct". Question: Photosynthesis is concentrated in the mesophyll of the leaf. The "mesophyll" is a type of "tissue" composed of two layers or arrangements of chlorenchyma cells: an upper pallisade layer of tightly packed parenchyma (called pallisade chlorenchyma), and a lower spongy layer of loosely packed parenchyma (or spongy chlorenchyma). One could regard these layers as different tissue types ("simple tissues") and the mesophyll as a "complex tissue". The mesophyll is packed between two protective layers of epidermal cells (tissue: epidermis of epidermal cells and cuticle), which along with the vascular tissue and perhaps other structural tissues form an "organ" called the leaf whose primary function is food production for the "organsim": the plant. Leaves usually have a stem-like structural part called a petiole by which they attach to another plant organ termed the stem. « Chapter 3 

Botany Both this Guide and all articles in "Wikipedia" are that can be added to or edited by anyone. It is an opportunity for the user of these documents to contribute information, or even state given information more clearly, simply by editing a page. As a student with a textbook and a lecturer (teacher), you may find yourself in possession of useful facts, another point of view on existing facts, or a report you prepared of exceptional quality. Any of these can be added to an appropriate page in this Guide or the "Wikipedia". However, this caution is strongly advised: "Do not place into the Guide any text or pictures taken verbatim (or close to verbatim) from a text book, web site, or other copyrighted source without permission of the copyright holder". In general, this means, anything you submit should be your own work. To learn how to edit or contribute material to this textbook, first read the introduction at: "". 

Gibbs free energy is represented by the following equation: The basic principle is that "total" entropy increases. This increase can be because of an increase in the entropy of the chemicals, "ΔS", or because the reaction has produced heat, increasing the entropy of the environment. The Gibbs free energy lets us calculate the total increase in entropy, including the effects on the environment, without needing to know anything about the environment. At "low" temperatures, "ΔG" is approximately "ΔH", and nature favours the reaction with lowest energy products, which release the most heat. This may "reduce" the entropy of the "system", but the "increase" in the entropy of the "environment" more than compensates. At "high" temperatures, "ΔG" is approximately "-TΔS", and nature favours the reaction with high energy products, which may actually absorb heat. This may "reduce" the entropy of the "environment", but the "increase" in the entropy of the "system" more than compensates. Either way, the Gibbs free energy "always" decreases. 

Latin Mottos. Latin External Links. The internet is one of the best media to obtain Latin resources. A few sites to get you started are listed here. ^ Latin ^ 

This Page is incomplete For explanations of terms used in these grammar tables, check the Glossary. Declension of Nouns. 5th Declension Masculine/Feminine (each word has a set gender): res. NOTES 4/5th declensions are modified 3rd declensions, thus behave similarily. 3rd declension is either M/F/N, 4th declension is either M/F/N and 5th declension is either M/F. So for 3rd, 4th, and 5th declension is of most importance to memorise the gender because the adjective will still need to agree (ie have the same) with the noun in both case, gender, and number. The Vocative only changes for the 2nd declension masculine singular (however not for the nouns that leave the -us suffix when in nominative). There are a few exceptions in the 1st declension which are not feminine. Such nouns are poet-a (1st declension masculine, so to agree you need to use -us on the adjective) and naut-a. There are a few second declension nouns with irregularities. It is of most importance that you memorise them. When memorising a noun's meaning, make sure that you also memorise any irregularities the noun has, the gender, and the declension. Without doing this you may have trouble translating. For example, 2nd declension masculines have -us in the nominative singular; however, 4th declension masculines have -us in the nominative singular, nominative plural and accusative plural. This may get you confused if you do not memorise the declension of each noun. Single Declension Theory. If you look at the above list of declensions, you may feel that you are going to be overcome if you have to memorize all of it. Memorization is indeed the key, but it will be easier than you think. Each word in Latin has three parts: the root, the stem vowel, and the ending. There are five vowels in Latin, so there are five stem vowels: A,O,I,U,E. Any word not really having a stem vowel naturally was given to I. If you study the declension patterns enough you will see that there are many similarities in the declensions. All accusative singulars end in "m" except neuters that sometimes still do. All accusative plurals end in 's' except neuters that always end in "a". All genitive plurals end in "um", be it "ium, rum, uum" or whatever else. If you read enough Latin you will actually find that authors would switch a declension of a word at will if it made the sentence clearer. Thus many words that are of the 4th declension were sometimes written as if 2nd, and 5th declension as if 3rd and vice versa. Some ending patterns that we use for one declension may also be used in another, again to make the sentence clearer. Thus "filiabus" would be used in any sentence where we want to make it clear that we are talking about the Daughters and specifically not the Sons. We also see examples of this in "animabus". Some think that the "IS" used in the first and second declensions was actually an abbreviation of "bus". Some students find the ablative difficult since it sometimes looks like the nominative singular, dative plural or neither. All you really need to do to get the hang of this is to know that in the plural, the ablative always looks like the dative. If there is no prepostion in front of it then it is probably a dative, unless it is being used in an ablative construction that would likely be apparent. If there is a preposition and it looks like a dative, just remember that no preposition takes a dative, only ever ablative and sometimes accusative. Dative plural plus preposition equals ablative. The ablative singular is really just the root plus the stem vowel with no ending to speak of, since the preposition tells you the grammar of the word. Some students also get confused by words that are the same in the nominative singular and plural. Don't worry about that; the verb will tell you which it is since the verb always agrees with the nominative. It is thought that originally there was only one declension, but during the task of applying it to every word in natural speech it was found that some words naturally changed the way the basic sounds of the original declension worked. Here is one idea of the original basic declension. When speaking in everyday conversation, Latin speakers would shorten the word in their pronunciation so long as it still made sense. If you go to New Orleans you will likely hear someone say 'prowly'. This is not a new dish at a restaurant or a new new code name for the police, it is actually the local pronunciation of the word "probably". It is ok to shorten this word because when it is shortened everyone can still understand it. The Latins did the same thing. It is known that "I" can change to "E" and vice versa, so we can see that pattern in the dative singular of all declensions. The "IS" of the genitive was kept in the 3rd declension, shortened in the 1st, 2nd, and 5th and slightly altered in the 4th. The 'add "M"' rule of the accusative singular is seen in the altering of the stem vowel in the 2nd and 3rd declensions. The "ES" of the nominative plural became "E" in 1st, "I" in 2nd, and "US" in 4th. The "RUM" of the genitive plural was shortened in the 3rd and 4th declensions. The accusative plural as a rule never changes, but keep in mind that neuter words followed a different pattern of using "A" for the accusative plural, and all neuters took the accusative singular or plural for the nominative of the same number. The "BUS" of the dative/ ablative plural was shortened to "IS" in the 1st and 2nd declensions but kept in the others. The 4th declension has a tendency to copy the "IBUS" of the 3rd in some authors but retain the "UBUS" in others. When studying Latin declensions you really should strive to memorize the patterns, but also look to see how they are similar to other patterns in the language as it will help you to remember. Also note that the "IS" of the genitive singular is what eventually gave us the word "HIS" and the 'add 's'" rule of English. Latin also gave us the 'add 's'" to make a plural by way of the accusative plural. If you study other inflexive languages related to Latin, such as Greek, you will notice even similarities across languages. 

The Genitive. The genitive case is a descriptive case. The genitive case describes the following features of the described noun: Quite simply, a word in the genitive case is translated with the preposition "of". Note that Latin does not have a separate form for the possessive genitive ("Marcus's dog" vs "The dog of Marcus"), as English does. A word in the genitive case showing possession can be translated either way. Exercise 1. Indicate the word in the genitive: Agreeing with the Adjectives. When adjectives are used to describe nouns in the genitive case, they must have the same case, number, and gender as the noun to which it refers. Example. It's that simple. The Dative. The dative case, also known as the indirect object case indicates: Latin does not distinguish between "to" or "for", though this is sometimes the case in English: Example 1. 'For' is the preposition indicating a dative. 'For' can be used in some other constructs. To determine whether it is dative, analyse the meaning of the sentence (see Example 3). Practice will enable you to quickly spot the case of a noun in the sentence without much effort. Example 2. "He gave the book to John"; "He gave to John the book"; or "He gave John the book". This demonstrates how English can use prepositions to change word order and even 'presume' a certain preposition exists that has been left out, giving a dative construct. Also, the dative is used only for a noun Exercise 2: Translate into English. Note that "placeo" requires the dative case, as opposed to the accusative case. Verbs such as this are denoted with "(+dat.)" or similar abbreviations. Roman Numerals. The Romans did not use the Hindu-Arabic numerals we use today. They used their own symbols and own numeric system. We still use Roman Numerals today. Note the declensions of the first three numbers. "Nullus" is the Latin equivalent of zero, for example: "nullam puellam in agro video" means "I see no girl in the field". 

Soluciones a los ejercicios. "Solutions to the exercises" Llena los espacios en blanco con la forma verbal correcta del verbo "ser". "Fill the blank spaces with the correct form of the ver "ser" (to be)." Llena los espacios en blanco con la forma verbal correcta del verbo "estar"."Fill the blank spaces with the correct form of the ver "estar" (to be)." Llena los espacios en blanco con la forma verbal correcta del verbo "ser" o "estar". "Fill the blank spaces with the correct form of the ver "ser" or "estar" (to be)." Enlace a los Ejercicios "Link to exercices" Enlace a la lección 1 "Link to the lesson 1." 

Haloalkanes are otherwise simple alkanes that contain one or more members of the halogen family. In practice, the halogens found in organic molecules are chlorine (Cl), bromine (Br), fluorine (F), and iodine (I). Some texts refer to this class of compounds as halogenoalkanes or alkyl halides. This text (and the chemical literature) will frequently use the terms haloalkane and alkyl halide interchangeably. "Note:" The X in R-X represents a generic halogen atom. =Preparation= Methods for preparation are found elsewhere in this text: =Properties= Naming Haloalkanes. Haloalkanes are named by adding a prefix to the name of the alkane from which they are derived. The prefix denotes the particular halogen used. F = "Fluoro-"&lt;br&gt; Cl = "Chloro-"&lt;br&gt; Br = "Bromo-"&lt;br&gt; I = "Iodo-"&lt;br&gt; If other substituents need to be named, all prefixes are still put in alphabetical order. When necessary, numbers identify substituent locations. Physical properties. R-X bond polarity: C—F &gt; C—Cl &gt; C—Br &gt; C—I The difference in electronegativity of the carbon-halogen bonds range from 1.5 in C-F to almost 0 in C-I. This means that the C-F bond is extremely polar, though not ionic, and the C-I bond is almost nonpolar. Physical appearance: Haloalkanes are colourless when pure. However bromo and iodo alkanes develop colour when exposed to light. Many volatile halogen compounds have sweet smell. Boiling point: Haloalkanes are generally liquids at room temperature. Haloalkanes generally have a boiling point that is higher than the alkane they are derived from. This is due to the increased molecular weight due to the large halogen atoms and the increased intermolecular forces due to the polar bonds, and the increasing polarizabilty of the halogen. For the same alkyl group, the boiling point of haloalkanes decreases in the order RI &gt; RBr &gt; RCl &gt; RF.This is due to the increase in van der Waals forces when the size and mass of the halogen atom increases. For isomeric haloalkanes, the boiling point decrease with increase in branching. But boiling points of dihalobenzenes are nearly same; however the para-isomers have higher melting points as it fits into the crystal lattice better when compared to ortho- and meta-isomers. Density: Haloalkanes are generally more dense than the alkane they are derived from and usually more dense than water. Density increases with the number of carbon and halogen atom. It also increases with the increase in mass of halogen atom. Solubility: The haloalkanes are only very slightly soluble in water, but dissolves in organic solvents. This is because for dissolving haloalkanes in water the strong hydrogen bonds present in the latter has to be broken. When dissolved in organic (non polar) solvents, the intermolecular attractions are almost same as that being broken. Bond Length: C—F &lt; C—Cl &lt; C—Br &lt; C—I Larger atoms means larger bond lengths, as the orbitals on the halogen is larger the heavier the halogen is. In F, the orbitals used to make the bonds is 2s and 2p, in Cl, it's 3s and 3p, in Br, 4s and 4p, and in I, 5s and 5p. The larger the principal quantum number, the bigger the orbital. This is somewhat offset by the larger effective nuclear charge, but not enough to reverse the order. Chemical properties. Bond strength: C—F &gt; C—Cl &gt; C—Br &gt; C—I The orbitals C uses to make bonds are 2s and 2p. The overlap integral is larger the closer the principal quantum number of the orbitals is, so the overlap is larger in the bonds to lighter halogens, making the bond formation energetically favorable. Bond reactivity: C—F &lt; C—Cl &lt; C—Br &lt; C—I Stronger bonds are more difficult to break, making them less reactive. In addition, the reactivity can also be determined by the stability of the corresponding anion formed in solution. One of the many trends on the periodic table states that the largest atoms are located on the bottom right corner, implying that iodine is the largest and fluorine being the smallest. When fluorine leaves as fluoride (if it does) in the reaction, it is not so stable compared to iodide. Because there are no resonance forms and inductive stabilizing effects on these individual atoms, the atoms must utilize their own inherent abilities to stabilize themselves. Iodide has the greatest surface area out of these four elements, which gives it the ability to better distribute its negative charge that it has obtained. Fluorine, having the least surface area, is much more difficult to stabilize. This is the reason why iodine is the best leaving group out of the four halogens discussed. =Reactions= Determination of Haloalkanes: A famous test used to determine if a compound is a haloalkane is the Beilstein test, in which the compound tested is burned in a loop of copper wire. The compound will burn green if it is a haloalkane. The numbers of fluorine, chlorine, bromine and iodine atoms present in each molecule can be determined using the sodium fusion reaction, in which the compound is subjected to the action of liquid sodium, an exceptionally strong reducing agent, which causes the formation of sodium halide salts. Qualitative analysis can be used to discover which halogens were present in the original compound; quantitative analysis is used to find the quantities. Substitution reactions of haloalkanes. R-X bonds are very commonly used throughout organic chemistry because their polar bonds make them reasonably reactive. In a substitution reaction, the halogen (X) is replaced by another substituent (Y). The alkyl group (R) is not changed. Substitutions involving haloalkanes involve a type of substition called Nucleophilic substitution, in which the substituent Y is a nucleophile. A nucleophile is an electron pair donor. The nucleophile replaces the halogen, an "electrophile", which becomes a leaving group. The leaving group is an electron pair acceptor. Nuclephilic substition reactions are abbreviated as SN reactions. Example: Suggest a reaction to produce the following molecule. Answer: OR "Any halogen could be used instead of Br" Reaction mechanisms. Nucleophilic substitution can occur in two different ways. SN2 involves a backside attack. SN1 involves a carbocation intermediate. SN2 mechanism SN1 mechanism Comparison of SN1 and SN2 mechanism. Stereochemistry: &lt;br&gt; SN2 - Configuration is inverted (i.e. R to S and vice-versa). &lt;br&gt; SN1 - Product is a mixture of inversion and retention of orientation because the carbocation can be attacked from either side. In theory the products formed are usually racemic due to the 50% chance of attack from the planar conformation. Interestingly, the amount of the inverted product is often up to 20% greater than the amount of product with the original orientation. Saul Winstein has proposed that this discrepancy occurs through the leaving group forming an ion pair with the substrate, which temporarily shields the carbocation from attack on the side with the leaving group. &lt;br&gt; Rate of reaction:&lt;br&gt; SN2 - Rate depends on concentrations of both the haloalkane and the nucleophile. SN2 reactions are fast.&lt;br&gt; SN1 - Rate depends only on the concentration of the haloalkane. The carbocation forms much slower than it reacts with other molecules. This makes SN1 reactions slow.&lt;br&gt; Role of solvent:&lt;br&gt; SN2 - Polar aprotic solvents favored. Examples: Acetone, THF (an ether), dimethyl sulfoxide, n,n-dimethylformamide, hexamethylphosphoramide (HMPA).&lt;br&gt; Nonpolar solvents will also work, such as carbon tetrachloride (CCl4)&lt;br&gt; Protic solvents are the worst type for SN2 reactions because they "cage," or solvate, the nucleophile, making it much less reactive.&lt;br&gt; SN1 - Polar protic solvents favored. Examples: H2O, Formic acid, methanol.&lt;br&gt; Aprotic solvents will work also, but protic solvents are better because they will stabilize the leaving group, which is usually negatively charged, by solvating it. Nonpolar solvents are the worst solvent for SN1 reactions because they do nothing to stabilize the carbocation intermediate.&lt;br&gt; Role of nucleophile:&lt;br&gt; SN2 - Good nucleophiles favored&lt;br&gt; SN1 - Any nucleophile will work (since it has no effect on reaction rate)&lt;br&gt; Carbocation stability:&lt;br&gt; 3° carbon - most stable = SN1 favored&lt;br&gt; 2° carbon - less stable = either could be favored&lt;br&gt; 1° carbon - seldom forms = SN2 favored&lt;br&gt; CH3+ - never forms = SN2 favored&lt;br&gt; The reason why the tertiary carbocation is most favored is due to the inductive effect. In the carbocation intermediate, there is a resulting formal charge of +1 on the carbon that possessed the haloalkane. The positive charge will attract the electrons available. Because this is tertiary, meaning that adjacent carbon atoms and substituents are available, it will provide the most electron-density to stabilize this charge. Example. Predict whether the following reactions will undergo SN2 or SN1 and tell why. 1: 2: 3: Answers:&lt;br&gt; 1) SN2. Good nucleophile, polar solvent.&lt;br&gt; 2) SN1. Tertiary carbon, polar solvent. Very slow reaction rate.&lt;br&gt; 3) SN2. Primary carbon, good nucleophile, nonpolar solvent.&lt;br&gt; Grignard reagents. Grignard reagents are created by reacting magnesium metal with a haloalkane. The magnesium atom gets between the alkyl group and the halogen atom with the general reaction as stated below: R-Br + Mg → R-Mg-Br Gringard reagents are very reactive and thus provide a means of organic synthesis from haloalkanes. For example, adding water gives the alcohol R-OH. Basic: R-X + Mg → R-Mg-X For example (X=Cl and R=CH3): CH3-Cl + Mg → CH3MgCl Elimination reactions. With alcoholic potassium hydroxide, haloalkanes lose H-X and form the corresponding alkene. Very strong bases such as KNH2/NH3 convert vic-dihalides (haloalkanes with two halogen atoms on adjacent carbons) into alkynes. 

The Arrival of Columbus. Christopher Columbus and three ships - the "Niña", the "Pinta", and the "Santa Maria" - set sail on August 3, 1492. On October 12, a lookout cried out that he had sighted land. The crew set foot on an island that day, naming it San Salvador. It is unknown which exact island was discovered by Columbus. (Note that the island presently called San Salvador is so-called in honor of Columbus' discovery; it is not necessarily the one on which Columbus set foot.) The Native Americans inhabiting the islands were described as "Indians" by Columbus, who had believed that he had discovered the East Indies (modern Indonesia). In reality, he had found an island in the Caribbean. He continued to explore the area, returning to Spain. Columbus' misconception that he found Asia was corrected years later by the Italian explorer Amerigo Vespucci, after whom "America" may be named. The Protestant Reformation. In Europe, the power of the Pope and the influence of Catholicism was undoubted. The Catholic religion affected every aspect of politics on the continent. However, in the sixteenth century, the conditions were ripe for reform. Gutenberg's printing press made the spread of ideas much easier. The influence of nationalism grew, and rulers began to resent the power possessed by the Pope. The Protestant movement may have commenced earlier, but the publication of "Ninety-Five Theses" by Martin Luther in 1517 spurred on the revolution within the Church. Luther attacked the Church's theology, which, he believed, misrepresented The Bible and placed too much authority in the hands of the clergy, and wished to reform the Church. After being excommunicated, Luther published many books on Reform. Luther's works were most influential in Germany and Scandinavia. Persons other than Luther championed the cause of Reform. In Switzerland, Huldreich Zwingli advanced Protestant ideas, which mostly affected his home country. Similarly, Frenchman John Calvin helped the spread of Protestantism in France and the Netherlands. English Protestantism resulted from the direct influence of the British monarch. Henry VIII (1509-1547) sought to divorce his wife, Catherine of Aragon, because she had failed to produce a viable male heir to the throne. When his divorce led to excommunication by the Pope, Henry simply declared the entire country free of Catholic domination and a bastion of Protestantism. Henry reasoned that England could survive under its own religious regulation (Anglican) and he named himself head of the church. Elizabethan England. Elizabethan Succession&lt;br&gt; After Henry VIII died, he was succeeded by his son Edward VI (1547-1553) who reigned briefly before dying. Edward's death led to the ascension of Henry's daughter by Catherine, Mary I (1553-1558). A staunch Catholic, Mary sought to return England back to the Catholic church. Her religious zeal and persecution of Protestants earned her the nickname, "Bloody Mary." After a short reign, she was succeeded by her half-sister, Elizabeth. Elizabeth I (1558-1603) was the daughter of Henry's second wife, Anne Boleyn. Her ascendency to the throne resulted when neither of her half siblings, Edward and Mary, produced an heir to the throne. Religious Reform&lt;br&gt; Under her siblings' reign, the nation constantly battled religious fervor as it sought to identify itself as either Protestant or Catholic. Henry VIII had severed ties with the Roman Catholic Church upon his excommunication after divorcing Catherine of Aragon. He established the Church of England (the precursor to the Anglican Church) as the official state religion and named himself, not the Pope, as its head. Under Mary, the country returned to Catholic rule. The Elizabethan Age brought stability to English government. Elizabeth sought a compromise (the Elizabethan Settlement) which returned England to a nation governed by Protestant theology with a Catholic ritual. Elizabeth called Parliament in 1559 to consider the Reformation Bill that re-established an independent Church of England and redefined the sacrament of communion. Parliament also approved the Act of Supremacy, establishing ecclesiastical authority with the monarch. Economic Reform&lt;br&gt; Elizabeth's far more important response was to stabilize the English economy following the 1551 collapse of the wool market. To respond to this economic crisis, Elizabeth used her power as monarch to shift the supply-demand curve. She expelled all non-English wool merchants from the empire. Her government placed quotas on the amount of wool that could be produced while also encouraging manors to return to agricultural production. She also started trading directly with the Spanish colonies in direct violation of their tariff regulations. This maritime violation would later result in an attack on England by the Spanish Armada in 1588. Queen Elizabeth was a very popular monarch. Her people followed her in war and peace. She remained unmarried until her death, probably through a reluctance to share any power and preferring a series of suitors. This gave her the name, "the Virgin Queen," and in honor of her, a colony was named "Virginia" a few years after her death. In the aftermath of the Armada's overwhelming defeat and building on the development of a strong fleet started by Henry VIII, England began to gain recognition as a great naval power. Nationalism in England increased tremendously. Thoughts of becoming a colonial power were inspired. These thoughts were aided by the fact that the defeated Spanish lost both money and morale, and would be easy to oppose in the New World. Early Colonial Ventures. Richard Hakluyt In 1584, Richard Hakluyt proposed a strong argument for expansion of English settlement into the new world. With his "Discourse Concerning Western Planting," Hakluyt argued that creating new world colonies would greatly benefit England. The colonies could easily produce raw materials that were unavailable in England. By establishing colonies, England would assure itself of a steady supply of materials that it currently purchased from other world powers. Second, inhabited colonies would provide a stable market for English manufactured goods. Finally, as the economic incentives were not enough, the colonies could provide a home for disavowed Englishmen. Roanoke The English had already begun the exploration of the New World prior to the Armada's defeat. In 1584, Queen Elizabeth granted Sir Walter Raleigh a charter authorizing him to explore the island of Roanoke, which is part of what is now North Carolina. Between 1584 and 1586, Raleigh financed expeditions to explore the island of Roanoke and determine if the conditions were proper for settlement. In 1586, about a hundred men were left on the island. They struggled to survive, being reduced to eating dogs. They were, however, rescued- except for fifteen men whose fate remained a mystery. After another expedition in 1587, another group of men, women, and children- a total of more than one-hundred people- remained on the island. Governor John White of the Roanoke colony discovered from a local Native American tribe that the fifteen men who were not rescued were killed by a rival tribe. While attempting to gain revenge, White's men killed members of a friendly tribe and not the members of the tribe that allegedly killed the fifteen men. Having thus strained relations with the Natives, the settlers could not survive easily. John White decided to return to England in 1587 and return with more supplies. When he returned, England faced war against Spain. Thus delayed, White could not return to Roanoke until 1590. When he did return, White discovered that Roanoke was abandoned. All that gave clue to the fate of the colony was the word "Croatan," the name of a nearby Native American tribe, carved out on to a tree. No attempt was made to discover the actual cause of the disappearance until several years later. There are only theories as to the cause of the loss of Roanoke. There are two major possibilities. Firstly, the settlers may have been killed by the Natives. Second, the settlers may have assimilated themselves into the Native tribes. But there is no evidence that settles the matter beyond doubt. Review Questions. Use the content in this chapter and/or from external sources to answer the following questions. Remember to properly cite any sources used. 

Patterns of Colonization. The islands of Great Britain changed greatly in the Renaissance, resulting in the Church of England, the British Civil War, and total transformation of economic, political, and legal systems. Yet through this time, despite pressure from other nations and America's own Natives, a diverse set of English colonies were planted and thrived. These new colonies were funded in three different ways. In one plan, corporate colonies were established by joint stock companies. A joint stock company was a project in which people would invest shares of stock into building a new colony. Depending on the success of the colony, each investor would receive profit based on the shares he had bought. This investment was less risky than starting a colony from scratch, and each investor influenced how the colony was run. These investors often elected their own public officials. (An example of a joint stock company on another continent was the British East India Company.) Virginia was settled in this way. Proprietary colonies were owned by a person or family who made laws and appointed officials as he or they pleased. Development was often a direct result of this ownership. Charles II granted William Penn the territory now known as Pennsylvania. Penn's new colony gave refuge to Quakers, a group of millennial Protestants who opposed the Church of England. (Quakers did not have ministers and did not hold to civil or religious inequality, making them a dangerous element in hierarchical societies.) Penn was an outspoken Quaker and had written many pamphlets defending the Quaker faith. He also invited settlers from other countries and other Protestant minorities, and even some Catholics. Finally, royal colonies were under the direct control of the King, who appointed a Royal Governor. The resulting settlement was not always identical to England. For example, England had broken with Catholicism during the reign of Henry the Eighth, and the Old Faith was seen not only as religious heresy but the prelude to domination by other countries. Yet Maryland's grant of toleration of Catholics was granted as a boon from the British Crown. In 1634, Lord Baltimore appointed George Calvert of England to settle a narrow strip of land north of Virginia and south of Pennsylvania as a Catholic colony via a royal charter. Fifteen years later, in 1649, he signed the Act of Toleration, which proclaimed religious freedom for its colonists. Despite the original charter, Protestants later became the majority faith. After Lord Baltimore's death several years later, Margaret Brent, the wife of an esteemed landowner in Maryland, executed his will as governor of the colony. She defied gender roles in the colonies by being the first woman of non-royal heritage to govern an English colony. Massachusetts Bay Colony. The "Massachusetts Bay Colony", another corporate colony, was founded as a place far from England where its religious dissenters could live. The Puritans, a group of radical Protestants who wanted what they called a return to the faith of the Bible, suffered torture and execution because they disagreed with the official Church of England. In 1620, forty-one Puritans (who called themselves Pilgrims) sailed for the new world. Their own contemporary accounts show that the Pilgrims originally intended to settle the Hudson River region near present day Long Island, New York. Once Cape Cod was sighted, they turned south to head for the Hudson River, but encountered treacherous seas and nearly shipwrecked. They then decided to return to Cape Cod rather than risk another attempt to head south. After weeks of scouting for a suitable settlement area, the Mayflower's passengers finally landed at Plymouth in present-day Massachusetts on December 26, 1620. They called it Massachusetts after the name of the Indian tribe then living there. William Bradford, who was selected as a governor after the death of John Carver, wrote a journal that helps us to better understand the hardships colonists endured, encounters with the Native Americans, and ultimately, the success of the colony. The Pilgrims agreed to govern themselves in the manner set forth in the Mayflower Compact, which signed on the Pilgrims' ship, The Mayflower. After two years they abandoned the communal form of partnership begun under the Compact and in 1623 assigned individual plots of land to each family to work. Ten years later, the joint-stock Massachusetts Bay Company acquired a charter from King Charles of England. The colony of Plymouth was eventually absorbed by Massachusetts Bay, but it remained separate until 1691. A large group of Pilgrims later migrated to the new colony of Massachusetts Bay. In keeping with its mother Church of England, the colony did not provide religious freedom. It only allowed (male) Puritans the right to vote, established Puritanism as the official religion of the colony in The Act of Toleration, and punished people who did not go to their Church. New York. Other countries used the joint-stock company to fund exploration. In 1609, the Dutch East India company discovered a territory on the eastern coast of North America, from latitude 38 to 45 degrees north. This was an expedition in the yacht Halve Maen ("Half Moon") commanded by Henry Hudson. Adriaen Block and Hendrick Christiaensz explored the territory from 1611 until 1614. In March of 1614 the States General, the governing body of the Netherlands, proclaimed exclusive patent for trade in the New World. The States General issued patents for development of New Netherland as a private commercial venture. Ft. Nassau was swiftly built in the area of present day Albany to defend river traffic and to trade with Native Americans. New Netherland became a province of the Dutch Republic in 1624. The northern border was then reduced to 42 degrees north, as the English had encroached north of Cape Cod. According to the Law of Nations, a claim on a territory required not only discovery and charting but settlement. In May 1624 the Dutch completed their claim by landing thirty Dutch families on Noten Eylant, modern Governors Island. In the next few decades incompetent directors-general ran New Netherland. The settlers were attacked by Native Americans, and British and Dutch conflicts seemed destined to destroy the colony. All that changed when Peter Stuyvesant was appointed Director-General in 1647. As he arrived he said, "I shall govern you as a father his children". He expanded the colony's borders. He oversaw conquest of the one settlement of northernmost Europe, New Sweden, in 1655. He resolved the border dispute with New England in 1650. He improved defenses against Native American raids, and the population of the colony went from 500 in 1640 to 9,000 by 1664. But in August of 1664, four English warships arrived in New York Harbor demanding the surrender of the colony. At first, Stuyvesant vowed to fight, but there was little ammunition and gunpowder. He received weak support from the overwhelmed colonists, and was forced to surrender. New Netherland was subsequently renamed "New York", in honor of the British Duke of York. In an attempt to gain supremacy over trade, the English waged war against the Dutch in 1664. The English took control over the Dutch harbor of New Amsterdam on the Atlantic coast of America. James, the brother of King Charles II, received the charter for New Amsterdam and the surrounding Dutch territory. In 1673 the Dutch, lead by Michiel de Ruyter, briefly reoccupied New Netherland again, this time naming it New Orange. After peace was made, ending the Third Anglo-Dutch War, they agreed to return it to the English. Patterns of Colonization in the Other Early Colonies. The territory of Carolina, named after the British King Charles I, was granted as a proprietary colony to eight different nobles. The proprietors divided Carolina into two separate colonies -- "North Carolina" and "South Carolina". Four colonies were formed by division from already extant larger territories. When New Holland was taken to become New York, King James granted a portion of the territory, present-day "New Jersey", to Lord Berkeley and Sir George Cartaret, while retaining present-day New York for himself as a proprietary colony. Sir George had come from the Isle of Jersey, and the new colony was named accordingly. Another portion of the territory became the crown colony Connecticut. This colony was also named for its native tribe of Indians. A corner of Pennsylvania which was not peopled by Quakers separated in 1704 to become the colony of "Delaware". This was given the name of Thomas West, Third Baron De La Warr, a nobleman under Queen Elizabeth and a noted adventurer. "Rhode Island" was a unique experiment in religious and political freedom. Massachusetts banished Roger Williams after he began asserting that Jesus Christ meant for the Church to be separate from the governing authority. This dissenter from the Church of England, and then from the Puritans, became the first American Baptist. After many adventures in other colonies, he bought land from the Narragansett Indians for a new settlement. Providence was meant to be a colony free from religious entanglements and a refuge for people of conscience. He was later followed by Anne Hutchinson. She had outraged Boston divines because she was a woman who preached, and because she believed that one's works were not always tied to grace, unlike the Puritans. She also bought land from the Indians. On this was the settlement subsequently named Portsmouth, and afterward a dissident sister town, Newport. The colony was partially based upon Aquidneck Island, later called Rhode Island for unknown reasons, and the entire establishment eventually took its name from that place. Georgia was another proprietary colony, named after King George I, with a charter granted to James Oglethorpe and others in 1732. It was intended as a "buffer" colony to protect the others from attacks from the Florida Spanish and the Louisiana French. Because of this, Georgia was the only colony to receive funds from the Crown from its founding. The laws in Great Britain put people in prison for debt. Many of these people were shipped from overcrowded jails to freedom in the wilds of Georgia colony. America was already seen as a land of prosperity, and Oglethorpe hoped that the ex-prisoners would soon become honest and rich. However, few of the prisoners of London jails knew how to survive in the new wilderness. Portrait of the British Colonies. The Colonies are often considered as three groups: New England (New Hampshire, Massachusetts, Rhode Island, Connecticut), the Southern Colonies (Maryland, Virginia, the Carolinas, and Georgia), and the Middle Colonies (New York, New Jersey, Pennsylvania and Delaware). Sometimes the Carolinas and Georgia are counted as separate from the Chesapeake Colonies. Each group had geographic and economic characteristics. New England's rocky soil only encouraged small farms, and its agricultural opportunities were limited. Thus it focused on fishing, forestry, shipping, and small industry to make money. Richer land in the Southern colonies was taken over by individual farmers who grasped acreage. This created large plantation farms that grew tobacco, and later cotton. Farms in the Carolinas also farmed sugar, rice, and indigo. In the 17th century, these were farmed by indentured servants, people who would work for a period of years in return for passage to America and land. Many of these servants died before their indentures ended. A group of indentured servants rose up in Bacon's Rebellion in 1676. After Bacon's Rebellion, plantations began using African slaves instead. Even after release from indenture, many of these white people remained in the economic lower classes, though not subject to the slave codes, which became more harsh as time passed, denying almost all liberty to slaves in the southern colonies. By the American Revolution, one in five colonists was an African slave. And the products produced by slavery in the South were consumed and traded by towns in the Middle Colonies and New England. Few people questioned the slave economy. The Middle Colonies had medium-sized farms. These colonies also had people from many different cultures with many different beliefs. Individuals in these states used indentured servants, and later slaves, but there was not the concentration of masses of slave labor found in the Southern colonies. Another distinction lies in religious practices. New England was mostly Congregationalist, with some admixture of Presbyterian congregations and the religious non-conformists who called themselves Baptists. These were all descendants of dissenters before and during the British Civil War. The South was mostly Anglican, cherishing religious and secular traditions and holidays. The Middle Colonies held small groups of people from Holland, German lands, and even Bohemia, and they brought a welter of Catholic and Protestant faiths. Among the whites sent to the colonies by English authorities were many Scots-Irish people from Ulster. These had been Calvinist Protestants in the middle of a Irish Catholic majority, at odds both with them and with England. This minority settled in the frontier region of the Appalachian Mountains and eventually beyond in the Ohio and Mississippi country. In America their desire for land and freedom pushed the colonial boundary westward at little cost to the government, and provided an armed buffer between the eastern settlements and Native American tribes which had been driven away from the seaboard. Colonial frontiersmen endured a very harsh life, building their towns and farms by hand in a dense wilderness amid economic deprivation and native attack. Each colony developed its own areas of edification and amusement, depending upon the local faith and the local capacities. The culture of the South recorded early interest in musical theater, with Charleston, South Carolina and Williamsburg, Virginia as hubs of musical activity. A performance of Richard III, the first professional production of Shakespeare in America, took place in New York City in 1750. And preachers, lecturers, and singers entertained the colonists. Their commonalities were stronger than their differences. All three regions shared a population mostly derived from the British Isles. All had terrible roads, and all had connections to the Atlantic Ocean as a means of transportation. And all were tied to the Atlantic economy. Atlantic merchants used ships to trade slaves, tobacco, rum, sugar, gold, silver, spices, fish, lumber, and manufactured goods between America, the West Indies, Europe and Africa. New York, Philadelphia, Boston, and Charleston were the largest cities and main ports at that time. Early Technology. The first wave of colonists used hand labor to cultivate their farms, and established such land-based crafts such as pottery and tanning. As later ships brought cattle and horses, draft animals became part of the economy. Indentured servants, and then slaves kidnapped from Africa, were imported. This was when larger plantations began to be founded. In the latter part of the eighteenth century small-scale machine-based manufacturing began to appear. Individuals started to dig for coal and iron ore. New England used the latter to begin making building tools and horseshoes. A new textile industry arose, dependent in part upon Southern cotton. Powered by wood or coal and fed by the need for strong metal, household forges pioneered new techniques of iron-making. The blacksmith and the tinsmith became part of large settlements. Colonies started making mechanized clocks, guns, and lead type for printing. Mercantilism, Salutary Neglect and British Interference. The American colonies, entirely new societies separated by an ocean from Great Britain, believed they had the right to govern themselves. This belief was encouraged by Great Britain's Glorious Revolution and 1689 Bill of Rights, which gave Parliament the ultimate authority in government. A policy of relatively lax controls or Salutary Neglect ended in increased British regulation resulting from the policy of mercantilism, and seen through the Lords of Trade and the later Navigation Acts. Mercantilism. Parliament placed controls on colonial trade in obedience to the economic policy of mercantilism. This was the idea that a nation's economic power depended on the value of its exports. A country could use its colonies to create finished goods, rather than raw materials. These could be traded to other countries, thus increasing the strength of the colonizing nation. This policy had been put forth by a Frenchman named Jean-Baptiste Colbert. It seemed tailor-made for Great Britain. Spain had American gold as its economic base, and France had American furs. England had neither of these. But it had American cotton, molasses, and tobacco, as well as its state-of-the-art ships. Prior to the mid-1700's, the colonies had enjoyed a long period of "salutary neglect", where the British largely let the colonies govern themselves. This now ended. The Lords of Trade. In an attempt to enforce mercantilism policies, King Charles II created the Lords of Trade as a new committee on the Privy Council. The Lords of Trade attempted to affect the government of the colonies in a manner beneficial to the English, rather than to the colonists. The Lords of Trade attempted to convert all American colonies to royal ones so that the Crown could gain more power. Under King James II, the successor to Charles II, New York, New Jersey, and the Puritan colonies were combined into the Dominion of New England in 1687. However, the Dominion only lasted a brief time. King James II, a Catholic, was seen as a threat by British Protestants. James was overthrown (he was technically held abdicated by Parliament) in the bloodless Glorious Revolution of 1688. In 1689, James' daughter Mary II and her husband William III took the throne as joint rulers. However, the British Parliament actually held the power. The Dominion of New England was dissolved, the various separate colonies were reestablished, and the Lords of Trade were abandoned (replaced by a Board of Trade, a purely advisory body). Navigation Acts. Beginning in 1660, the Parliament of England passed the Navigation Acts to increase its benefit from its colonies. The Acts required that any colonial imports or exports travel only on ships registered in England, meaning that only England could have the shipping power and the fees derived from them. They forbid the colonies to export tobacco and sugar to any nation other than England. (Tobacco was then used as medicine, and sugar was used to make alcohol, also a medicine.) And the colonies could not import anything manufactured outside England unless the goods were first taken to England, where taxes were paid, and then to the colonies. In the 1730s, The Sugar Act established a tax of six pence per gallon of sugar or molasses imported into the colonies. By 1750, Parliament had begun to ban, restrict, or tax several more products. It tried to curtail all manufacture in the colonies. This provoked much anger among the colonists, despite the fact that their tax burdens were quite low, when compared to most subjects of European monarchies of the same period. Colonists hated the Navigation Acts because they believed they would be more prosperous and rich if they could trade on their own behalf. They also believed that some vital resources would not be found in Britain. Indians in the 1700s. Indians of the Great Plains: Today, the area where the Indians of all the Great Plains lived is located from the Rocky mountains to the Mississippi River. During the 1700s, there were about 30 tribes that lived on the Great Plains. These tribes tended to rely on buffalo as their food source as well as other daily needs, such as clothing. Not only did Indians, specifically women, make their clothing out of buffalo, but also out of deer. Women would soak the deer or buffalo and scrape off the hair of the dead animal. Also, Indian tribes traded with one another. The number of horses an individual owned was a sign of wealth; Indians would trade their horses for food, tools, weapons(such as guns), and hides. Since the tribes spoke many different languages from one another, they had to use sign language to be able to trade with each other. Philadelphia Election Riot. A riot broke out on election day in Philadelphia in 1742 as a result of the Anglican population disagreeing with the Quaker majority. The riot stemmed over a power struggle between the Anglican and Quaker population. The Quakers had a history of political dominance over Philadelphia. The German population backed the Quaker vote because of the Quaker Pacifism which would protect from higher taxes and ultimately the draft. On election day, the Anglicans and sailors fought with the Quakers and Germans. The Quakers were able to seek shelter in the courthouse and complete the election. The Anglican party lost the election and 54 sailors were jailed following the riot. Education. As the three sections of the colonies through the 1700s were made up of people with different interests, they provided differing sorts of education for their children. Although there were commonalities -- a rich family in any of the three regions might send a son to Europe for his education -- people in different colonies tended to educate in differing ways. New England's motives for education were both civil and religious. The good citizen had to know his or her Bible. The Massachusetts General School Law of 1647 stated that if more than 50 families lived in a community, a schoolteacher must be hired. This law gave a justification: "It being one chief project of that old deluder, Satan, to keep men from the knowledge of the Scriptures, as in former times by keeping them in an unknown tongue, so in these latter times by persuading from the use of tongues, that so that at least the true sense and meaning of the original might be clouded and corrupted with love and false glosses of saint-seeming deceivers; and to the end that learning may not be buried in the grave of our forefathers, in church and commonwealth, the Lord assisting our endeavors." This was the Pilgrim ethos, set up in opposition to what they saw as the ignorance imposed by tyrants. Both boys and girls were often taught to read the Bible by their parents, perhaps with the aid of a horn book, an alphabet and syllabary page covered by a protective layer of horn. In addition to being able to read the Bible, a Christian ought to be able to govern in his society. ("His" society: for government was the province of godly, property-holding men, rather than women.) To obtain this youths had to gain a classical education -- that is, one based thoroughly on Latin. The 1647 law was the beginning of the American grammar school, which initially taught Latin, but later included practical subjects such as navigation, engineering, bookkeeping, and foreign languages. Most of the schools opened in the colonial era were private. However, they had been preceded by the first public-supported school, the Boston Latin School, in 1635. It had a rigorous education, and as a result, few students. Harvard was the first university in America, founded in 1636 and originally intended to teach Protestant clergy. Because of the small number of people graduating from the classical curriculum, attendance was low. Some people jumped directly from the classical curriculum to the University, sometimes entering Harvard as young as 14 or 15 years old. Cotton Mather graduated Harvard at 15, an exception only because of his extreme precocity. In private schools, boys and girls learned penmanship, basic Math, and reading and writing English. These fed the various trades, where older children were apprenticed. Girls who did not become servants were often trained for domestic life, learning needlework, cooking, and the several days-long task of cleaning clothes. Like New England, the Middle Colonies had private schools which educated children in reading and writing. However, the basics were rarer. The further west one lived, the less likely one was to be able to go to school, or to read and write at all. Ethnic and religious sub-groups would have their own private schools, which taught their own children their own folk-ways. In none of the colonies was higher education certain. Secondary schools were very rare outside of such major towns as Boston, New York, Philadelphia, and Charleston. The Chesapeake experience was different again. Children could only could only read and write if their parents could. And the South had few schools, of any kind, until the Revolutionary era. Children in wealthy families would study with private tutors. Though wealthy girls might learn 'the womanly arts,' they would not have the same curriculum as their brothers. Martha Washington's granddaughter Eliza Custis was laughed at by her stepfather when "[I] thought it hard they would not teach me Greek and Latin because I was a girl -- they laughed and said women ought not to know those things, and mending, writing, Arithmetic, and Music was all I could be permitted to acquire." Middle class children might learn to read from their parents, and many poor children, as well as all black children, went unschooled. The literacy rates were lower in the South than the North until about the 19th century. In 1693 the College of William &amp; Mary was founded, Virginia's first University. As the 18th century wore on, it specialized not in theology for clergymen but in law. In 1701, the Collegiate College was founded. In 1718 it received funds from a Welsh governor of the British East India Company, Elihu Yale, and was renamed Yale College. These were later joined by several other universities, including Princeton in 1747. In the 18th century, astronomy, physics, modern history and politics took a bigger place in the college curriculum. Some colleges experimented with admitting Native American students in the 18th century, though not African-Americans. In 1640, "The whole Booke of Psalms Faithfully Translated into English Metre", commonly known as the Bay Psalm Book, was printed in Cambridge, Massachusetts. It was the first book written in the new world. The Bay Psalm Book was the first metrical English translation of the Biblical psalms. This famous and influential songbook was succeeded by a whole New England publishing industry. Sometime after 1687 the first "New England Primer" was published as an aid to childhood reading and spelling. An alternative to the classical curriculum emerged in Benjamin Franklin's American Academy, founded in Philadelphia in 1751. This body represented something closer to the modern American high school, offering vocational education. This sort of school later outnumbered the classical secondary school. However, Franklin's Academy was private as well, making such learning open only to those who could afford it. During this period colonists attempted to convert Native Americans to Christianity. Review Questions. 1. Choose one of the following colonies: New York, Virginia, Massachusetts, Georgia. In which of the three areas is it located? Why and how was it initially colonized? How did its immigrants and the religions they adhered to affect it? 2. Why did the British interfere with the colonies? 

The French and Indian War. "(The following text is from Wikipedia)" The French and Indian War (1754–1763) was the North American chapter of the Seven Years' War. The name refers to the two main enemies of the British, the royal French forces and the various American Indian forces allied with them. This conflict, the fourth such colonial war between the kingdoms of France and Great Britain, resulted in the British conquest of all of New France east of the Mississippi River, as well as Spanish Florida. France ceded control of French Louisiana west of the Mississippi to its Spanish ally, to compensate it for its loss of Florida. By the end of this war France kept only the tiny islands of Saint Pierre and Miquelon north of the Caribbean. These colonies today still allow France access to the Grand Banks. In Great Britain and France, the North American theatre of the Seven Years' War war usually has no special name, and so the entire worldwide conflict is known as the "Seven Years' War" (or the "Guerre de sept ans"). The "Seven Years" refers to events in Europe, from the official declaration of war in 1756 to the signing of the peace treaty in 1763. These dates do not correspond with the actual fighting in North America, where the fighting between the two colonial powers was largely concluded in six years, from the Jumonville Glen skirmish in 1754 to the capture of Montreal in 1760. Elsewhere the conflict is known by several names. In British North America, wars were often named after the sitting British monarch, such as King William's War or Queen Anne's War. Because there had already been a King George's War in the 1740s, British colonists named the second war in King George's reign after their opponents, and thus it became known as the "French and Indian War". This traditional name remains standard in the United States, although it obscures the fact that American Indians fought on both sides of the conflict. American historians generally use the traditional name or the European title (the Seven Years' War), and have also invented other, less frequently used names for the war, including the "Fourth Intercolonial War" and the "Great War for the Empire". Canadian francophones and English speakers both refer to it as the Seven Years' War ("Guerre de Sept Ans") or the War of the Conquest ("Guerre de la Conquête"), as the war in which New France was conquered by the British and became part of the British Empire. This war was also known as the "Forgotten War". Reasons for war. The French and Indian War began less than a decade after France and Great Britain had fought on opposing sides in the European War of the Austrian Succession (1740–1748). One cause for the conflict was territorial expansion. Newfoundland's Grand Banks were fertile fishing grounds and coveted by both sides. Both sides also wanted to expand their territories for trapping furs to trade, and for other pursuits that aided their economic interests. Both the British and the French used trading posts and forts to claim the Ohio Country, the vast territory between the Appalachian Mountains and the Mississippi River, from the Great Lakes to the Gulf of Mexico. English claims resulted from royal grants with no definite western boundaries. La Salle had claimed the Mississippi River for France: its drainage area includes the Ohio River Valley. Both Great Britain and France took advantage of Native American factions to secure these claims, to protect their territories, and to keep the other from growing too strong. A second cause was political &amp; religious ideology. The English Protestant colonists feared papal influence in North America. New France was administered by French governors and Roman Catholic hierarchy. French missionaries included Armand de La Richardie. English history was told as a story of freedom from Catholic (i.e., foreign) influence. French control over North America represented a threat to Great Britain. In their turn, the French feared English anti-Catholicism, in a time when Catholics were still being persecuted under English law. Declaration and Action Anticipating the War. Céloron's expedition. In June 1747 the Governor-General of New France, the Marquis de la Jonquière, ordered Pierre-Joseph Céloron to lead an expedition to the Ohio Country to remove British influence from the area. Céloron was also to confirm allegiance of the Native Americans in the Ohio territory to the French crown. Céloron's expedition consisted of 213 soldiers of the Troupes de la marine (French Marines) transported by twenty-three canoes. The expedition left Lachine on June 15, 1749, and two days later reached Fort Frontenac. It then continued along the shoreline of present-day Lake Erie. At Chautauqua Portage (Barcelona, New York), it moved inland to the Allegheny River. The troop headed south to the Ohio River at present-day Pittsburgh, where Céloron buried lead plates engraved with the French claim to the Ohio Country. Whenever British merchants or fur-traders were encountered by the French, they were informed of the illegality of being on French territory and told to leave the Ohio Country. When the expedition arrived at Logstown, the Native Americans there informed Céloron that they owned the Ohio Country, and they would trade with the British, despite anything the French said. Céloron continued the expedition. At its farthest point south, it reached the junction between the Ohio River and the Miami River, just south of the village of Pickawillany. Here lived the old Chief of the Miami tribe, whom Céloron called "Old Britain." When Céloron arrived at Pickawillany, he informed the elderly Chief of "dire consequences" of continuing to trade with the British. "Old Britain" ignored the warning. After this meeting, Céloron and his expedition began the trip home, reaching Montreal only on November 10, 1749. In his report, Céloron wrote: "All I can say is that the Natives of these localities are very badly disposed towards the French, and are entirely devoted to the English. I don't know in what way they could be brought back." Langlade's expedition. On March 17, 1752, Governor-General de la Jonquière died. His temporary replacement was Charles le Moyne de Longueuil. It was not until July 1, 1752 that Ange Duquense de Menneville arrived in New France to take over the post. In the spring of 1752, Longueuil dispatched an expedition to the Ohio River area. The expedition was led by Charles Michel de Langlade, an officer in the Troupes de la marine. Langlade was given 300 men, some French-Canadians, and others members of the Ottawa tribe. His objective was to punish the Miami of Pickawillany for continuing to trade with the British. At dawn on June 21, 1752, the war party attacked the British trading center at Pickawillany, killing fourteen people of the Miami nation, including "Old Britain." The expedition then returned home. Marin's expedition. In the spring of 1754, Paul Marin de la Malgue was given command of a 2,000 man force of Troupes de la Marine and Aboriginals. His orders were to protect the Ohio from the British. Marin followed the route that Céloron had mapped out four years before. However, where Céloron had buried lead plates, Marin was constructing and garrisoning forts. The first fort that was constructed by Paul Marin was at Presque Isle (Erie, Pennsylvania) on Lake Erie's south shore. He then had a road built to the headwaters of Rivière aux Boeuf (now known as Waterford, Pennsylvania). Marin then constructed a second fort at Le Boeuf, designed to guard the headwaters of the Rivière aux Boeuf. Tanaghrisson's proclamation. On September 3, 1753, Tanaghrisson (d. 1754), Chief of the Mingo, arrived at Fort Le Boeuf. One tradition states that Tanaghrisson hated the French because they had killed and eaten his father. Tanaghrisson told Marin, "I shall strike . . .", threatening the French. The show of force by the French had alarmed the Iroquois in the area. They sent Mohawk runners to William Johnson's manor in Upper New York. Johnson, known to the Iroquois as "Warraghiggey", meaning "He who does big business", had become a respected member of the Iroquois Confederacy in the area. In 1746, Johnson was made a colonel of the Iroquois, and later a colonel of the Western New York Militia. At Albany, New York, there was a meeting between Governor Clinton of New York and Chief Hendrick, as well as other officials from a handful of American colonies. Chief Hendrick insisted that the British abide by their obligations and block French expansion. When an unsatisfactory response was offered by Clinton, Chief Hendrick proclaimed that the "Covenant Chain", a long-standing friendly relationship between the Iroquois Confederacy and the British Crown, was broken. Dinwiddie's reaction. Governor Robert Dinwiddie of Virginia found himself in a predicament. Many merchants had invested heavily in fur trading in Ohio. If the French made good on their claim to the Ohio Country and drove out the British, then the Virginian merchants would lose a lot of money. Dinwiddie could not possibly allow the loss of the Ohio Country to France. In October 1753 he wrote a letter to the commander of the French forces in the Ohio Country, Jacques Legardeur de Saint-Pierre, demanding an immediate French withdrawal. To deliver it he delegated Major "George Washington" of the Virginia militia. Major Washington left for Fort Le Boeuf on the 31st of October, along with his interpreter Jacob Van Braam and several other men. A few days later, Washington and his party arrived at Wills Creek (Cumberland, Maryland). Here Washington enlisted the help of Christopher Gist, a surveyor who was familiar with the area. They arrived at Logstown on November 24, 1753. At Logstown, Washington met with Tanaghrisson, who was angry over the French military encroachment upon his land. Washington convinced Tanaghrisson to accompany his small group to Fort Le Boeuf. On December 12, 1753, Washington and his men reached Fort Le Boeuf. Commander Saint-Pierre invited Washington to dine with him that evening. Over dinner, Washington presented Saint-Pierre with the letter from Dinwiddie. Saint-Pierre was civil in his response, saying, "As to the Summons you send me to retire, I do not think myself obliged to obey it." The French explained to Washington that France's claim to the region was superior to that of the British, as René-Robert Cavelier, Sieur de a Salle (1643–1687) had explored the Ohio Country nearly a century earlier. Washington's party left Fort Le Boeuf early on December 16, 1753. By January 16, 1754, they had arrived in Williamsburg, Virginia. In his report, Washington stated, "The French had swept south." They had constructed and garrisoned forts at Presque Isle, Le Boeuf and Venango. War. The French and Indian War was the last of four major colonial wars between the British, the French, and their Native American allies. Unlike the previous three wars, the French and Indian War began on North American soil and then spread to Europe, where Britain officially declared war on France on May 15, 1756, marking the beginnings of the Seven Years' War in Europe. Native Americans fought for both sides, but primarily alongside the French (with one exception being the Iroquois Confederacy, which sided with the American colonies and Britain). The first major event of the war was in 1754. Lieutenant Colonel George Washington, then twenty-one years of age, was sent to negotiate boundaries with the French, who did not give up their forts. Washington led a group of Virginian (colonial) troops to confront the French at Fort Duquesne (present day Pittsburgh). Washington discovered the French troops at the Battle of Jumonville Glen (about six miles or ten kilometers North-West of soon-to-be-established Fort Necessity). In the ensuing skirmish, a French Officer, Joseph Coulon de Jumonville, was killed. Washington pulled back a few miles and established Fort Necessity. The French forced Washington and his men to retreat. Meanwhile, the Albany Congress was taking place as means to discuss further action. Edward Braddock led a campaign against the French at Fort Duquesne in 1755. Washington was again among the British and colonial troops. Braddock employed European tactics -- bold, linear marches and firing formations -- and employed heavy cannon. This led to disaster at the Monongahela. The French and natives were heavily outmanned and outgunned. But they used superior tactics, taking cover behind trees and bushes to gun down and rout the British. Braddock was killed. Despite four close calls, Washington escaped unharmed and led the survivors in retreat. This stunning British defeat heralded a string of major French victories over the next few years, at Fort Oswego, Fort William Henry, Fort Duquesne, and Carillon, where veteran Montcalm famously defeated five times his number. The sole British successes in the early years of the war came in 1755, at the Battle of Lake George, which secured the Hudson Valley; and in the taking of Fort Beauséjour (which protected the Nova Scotia frontier) by Lieutenant Colonel Robert Monckton. A consequence of this last battle was the subsequent forced deportation of the Acadian population of Nova Scotia and the Beaubassin region of Acadia. In 1756 William Pitt became Secretary of State of Great Britain. His leadership, and France's continued neglect of the North-American theater, eventually turned the tide in favor of the British. The French were driven from many frontier posts such as Fort Niagara, and the key Fortress Louisbourg fell to the British in 1758. In 1759, the Battle of the Plains of Abraham gave Quebec City to the British, who had to withstand a siege there after the Battle of Sainte-Foy a year later. In September of 1760, Pierre François de Rigaud, Marquis de Vaudreuil-Cavagnal, the King's Governor of New France, negotiated a surrender with British General Jeffrey Amherst. General Amherst granted Vaudreuil's request that any French residents who chose to remain in the colony would be given freedom to continue worshiping in their Roman Catholic tradition, continued ownership of their property, and the right to remain undisturbed in their homes. The British provided medical treatment for the sick and wounded French soldiers, and French regular troops were returned to France aboard British ships with an agreement that they were not to serve again in the present war. Summary of the war in America In 1752 the French and their Native allies raided a trading outpost sited at modern day Cleveland rid the area of Pennsylvanians. In 1754 General George Washington attacked French soldiers and then became trapped in his poorly built, Fort Necessity in Great Meadows Pennsylvania and more than one-third of Washingtons's men shortly became casualties. Twenty-two year old Washington and his men surrendered and were allowed to leave back to Virginia. In July 1755, a few miles south of Fort Duquensne in Pennsylvania, the combined forces of French and Natives attacked British colonial troops that were preparing a to assault the fort. The aftermath that ensued would result in a British defeat and General Edward Braddock would be killed. Once London heard of this Britain declared war upon France and formally began the seven years war. After this the British feared that France would attempt to retake Nova Scotia and that the 12,000 French Nova Scotians would break their neutrality, so the British military forced around seven thousand French Nova Scotians from their homeland. This was history's first large-scale modern deportation, the French would be sent from Louisiana to the Caribbean and families would become torn apart. In July of 1758 The British had recaptured the fort at Loiusberg winning control of the St. Lawerence River. This would cut the major French supply route and open up more supply lines for the British. In the fall of 1758 the Shawnee and Delaware Natives accepted peace offerings from the British and the French abandoned Fort Duquesne. A decisive attack would happen in the fall of 1759 when General James Wolfe's forces defeated the French on the Plains of Abraham and thus taking Quebec. A year after this event the British would capture Montreal and the North American stage of the war would be over. Outcome. Though most of the North American fighting ended on September 8, 1760, when the Marquis de Vaudreuil surrendered Montreal — and effectively all of Canada — to Britain (one notable late battle allowed the capture of Spanish Havana by British and colonial forces in 1762), the war officially ended with the signing of the Treaty of Paris on February 10, 1763. The treaty sealed France's loss of all its North American possessions east of the Mississippi except for Saint Pierre and Miquelon islands off Newfoundland. All of Canada was ceded to Britain. France regained the Caribbean islands of Guadeloupe and Martinique, which had been occupied by the British. The economic value of these islands to France was greater than that of Canada at the time, because of their rich sugar crops; and the islands were easier to defend. However, the British were happy to take New France: defense was not an issue, and they had many sources of sugar. Spain gained Louisiana, including New Orleans, in compensation for its loss of Florida to the British. French Canada contained approximately 65,000 French-speaking Roman Catholic residents. Early in the war, in 1755, the British had expelled French settlers from Acadia. (Some of these eventually fled to Louisiana, creating the Cajun population.) Now at peace, and eager to secure control of its hard-won colony, Great Britain made concessions to its newly conquered subjects with the Quebec Act of 1774. The history of the Seven Years' War, particularly the siege of Québec and the death of British Brigadier General James Wolfe, generated a vast number of ballads, broadsides, images, maps and other printed materials, which testify to how this event continued to capture the imagination of the British public long after Wolfe's death in 1759. The European theatre of the war was settled by the Treaty of Hubertusburg on February 15, 1763. The war changed economic, political, and social relations between Britain and its colonies. It plunged Britain into debt, which the Crown chose to pay off with tax money from its colonies. These taxes contributed to the beginning the American Revolutionary War. Proclamation of 1763. "(The following text is taken from the Wikipedia article)" The Royal Proclamation of 1763 was issued October 7, 1763 by George III following Great Britain's acquisition of French territory in North America after the end of the Seven Years' War. The purpose of the proclamation was to make sure Britain could control its new territory for its The Proclamation in essence forbade Americans from settling or buying land west of the Appalachians. Colonists were angry because many already had land in that area. Additionally, the Proclamation gave the Crown a monopoly in land bought from Native Americans. Native land. In the fall of 1763, a royal decree was issued that prohibited the North American colonists from establishing or maintaining settlements west of an imaginary line running down the crest of the Appalachian Mountains. The proclamation acknowledged that Native Americans owned the lands on which they were then residing and white settlers in the area were to be removed. However, provision was made to allow specially licensed individuals and entities to operate fur trading ventures in the proscribed area. There were two motivations for this policy: To avoid warfare with the Indians. Neither side evidenced any affection for the tribes, but Indian conflicts were very expensive, and the British hadn't yet deployed enough soldiers in the West to keep the peace. Some Indians welcomed this policy, believing that separation from the colonies would allow them to resume their traditions. Others realized that the proclamation would, at best, only provide breathing room before the next onslaught of settlers. To concentrate colonial settlements on the seaboard where they could be active parts of the British mercantile system. British trade officials took it as a first priority to populate the recently secured areas of Canada and Florida (referring to the Treaty of Paris), where colonists could reasonably be expected to trade with the mother country. Settlers living west of the Appalachians would be highly self-sufficient and have little opportunity to trade with English merchants. The reaction of colonial land speculators and frontiersmen was immediate and negative. They believed their fight in the recent war had been "rewarded" by the creation of a vast restricted native reserve in the lands they coveted. Most concluded that the proclamation was only a temporary measure: a number ignored it entirely and moved into the prohibited area. Almost from its inception, the proclamation was modified to suit the needs of influential people with interests in the American West, both high British officials and colonial leaders. Beginning in 1764, portions of the Proclamation Line were adjusted westward to accommodate speculative interests. Later, in 1768, the first Treaty of Fort Stanwix formally recognized the surrender of transmontane lands claimed by the Iroquois. The Proclamation of 1763 was a well-intentioned measure. Pontiac’s Rebellion had inflicted a terrible toll on the frontier settlements in North America and the British government acted prudently by attempting to avoid such conflict in the foreseeable future. The colonists, however, were not appreciative and regarded the new policy as an infringement of their basic rights. The fact that western expansion was halted at roughly the same time that other restrictive measures were being implemented, made the colonists increasingly suspicious Almost immediately, many British colonists and land speculators objected to the proclamation boundary, since there were already many settlements beyond the line (some of which had been temporarily evacuated during Pontiac's War), as well as many existing land claims yet to be settled. Indeed, the proclamation itself called for lands to be granted to British soldiers who had served in the Seven Years' War. Prominent American colonists joined with land speculators in Britain to lobby the government to move the line further west. As a result, the boundary line was adjusted in a series of treaties with Native Americans. The Treaty of Fort Stanwix and the Treaty of Hard Labor (both 1768) and the Treaty of Lochaber (1770) opened much of what is now West Virginia and Kentucky to British settlement. Organization of new colonies. Besides regulating colonial expansion, the proclamation dealt with the management of newly ceded French colonies. It established government for four areas: Quebec, West Florida, East Florida, and Grenada. All of these were granted the ability to elect general assemblies under a royally appointed governor or a high council, which could then create laws and ordinances specific to the area in agreement with British and colonial laws. In the meantime, the new colonies enjoyed the same rights as native-born Englishmen, something that British colonists had been fighting over for years. An even bigger affront to the British colonies was the establishment of both civil and criminal courts complete with the right to appeal--but those charged with violating the Stamp or Sugar Act were to be tried in admiralty court, where the defendant was considered guilty until he or she could prove his or her innocence. Legacy. The influence of the Royal Proclamation of 1763 on the coming of the American Revolutionary War (1775–1783) has been variously interpreted. Many historians argue that the proclamation ceased to be a major source of tension after 1768, since the aforementioned treaties opened up extensive lands for settlement. Others have argued that colonial resentment of the proclamation contributed to the growing divide between the colonies and the Mother Country. In the United States, the Royal Proclamation of 1763 ended with the American Revolutionary War, because Great Britain ceded the land in question to the United States in the Treaty of Paris (1783). Afterwards, the U.S. government also faced difficulties in preventing frontier violence, and eventually adopted policies similar to those of the Royal Proclamation. The first in a series of Indian Intercourse Acts was passed in 1790, prohibiting unregulated trade and travel in Native American lands. Additionally, the U.S. Supreme Court case Johnson v. M'Intosh (1823) established that only the U.S. government, and not private individuals, could purchase land from Native Americans. The Royal Proclamation continued to govern the cession of aboriginal land in British North America, especially Upper Canada and Rupert's Land. The proclamation forms the basis of land claims of aboriginal peoples in Canada – First Nations, Inuit, and Metis. The Royal Proclamation of 1763 is thus mentioned in Section Twenty-five of the Canadian Charter of Rights and Freedoms. The Stamp Act and other Laws. In 1764, George Grenville became the British Chancellor of the Exchequer (minister of finance). He allowed customs officers to obtain general writs of assistance, which allowed officers to search random houses for smuggled goods. Grenville thought that if profits from smuggled goods could be directed towards Britain, the money could help pay off debts. Colonists were horrified that they could be searched without warrant at any given moment. Also in 1764, with persuasion from Grenville, Parliament began to impose several taxes on the colonists. The Sugar Act of 1764 reduced the taxes imposed by the Molasses Act, but at the same time strengthened the collection of the taxes. It also provided that British judges, and not juries, would try cases involving that Act. The next year, Parliament passed the Quartering Act, which required the colonies to provide room and board for British soldiers stationed in North America; the soldiers would serve various purposes, chiefly to enforce the previously passed acts of Parliament. Following the Quartering Act, Parliament passed one of the most infamous pieces of legislation: the Stamp Act. Previously, Parliament imposed only external taxes on imports. But the Stamp Act provided the first internal tax on the colonists, requiring that a tax stamp be applied to books, newspapers, pamphlets, legal documents, playing cards, and dice. The legislature of Massachusetts requested a conference on the Stamp Act; the Stamp Act Congress met in October that year, petitioning the King and Parliament to repeal the act before it went into effect at the end of the month, crying "taxation without representation." The act faced vehement opposition throughout the colonies. Merchants threatened to boycott British products. Thousands of New Yorkers rioted near the location where the stamps were stored. In Boston, the Sons of Liberty, a violent group led by radical statesman Samuel Adams, destroyed the home of Lieutenant Governor Thomas Hutchinson. Parliament did indeed repeal the Stamp Act, but additionally passed the Declaratory Act, which stated that Great Britain retained the power to tax the colonists, even without substantive representation. Believing that the colonists only objected to internal taxes, Chancellor of the Exchequer Charles Townshend proposed bills that would later become the Townshend Acts. The Acts, passed in 1767, taxed imports of tea, glass, paint, lead, and even paper. The colonial merchants again threatened to boycott the taxed products, reducing the profits of British merchants, who in turn petitioned Parliament to repeal the Townshend Acts. Parliament eventually agreed to repeal much of the Townshend legislation. But Parliament refused to remove the tax on tea, implying that the British retained the authority to tax the colonies despite a lack of representation. In 1773, Parliament passed the Tea Act, which exempted the British East India Company from the Townshend taxes. Thus, the East India Company gained a great advantage over other companies when selling tea in the colonies. The colonists who resented the advantages given to British companies dumped British tea overboard in the Boston Tea Party in December of 1773. &lt;br&gt;"The Boston Tea Party" In retaliation for the Boston Tea Party, Parliament passed the Coercive Acts, which were in the colonies known as the Intolerable Acts. Parliament reduced the power of the Massachusetts legislature and closed the port of Boston. Also, the Quartering Act was extended to require private individuals to lodge soldiers. Furthermore, Parliament allowed royal officials accused of crimes to be tried by a British, rather than a colonial, jury. First Continental Congress. In order to debate a response to the Intolerable Acts, all American colonies except for Georgia sent delegates to the First Continental Congress at Philadelphia. The Congress met in September 1774 and issued a Declaration of Rights and Grievances. When the Congress adjourned, it stipulated another Congress would meet if King George III did not meet the demands of the Declaration. When the Second Congress did meet, the military hostilities of the Revolutionary War had already begun, and the issue of Independence, rather than a redress of grievances, dominated the debates. Education. Literacy grew for both men and women during the 18th century. In New England and the Middle States, more middle-class girls were sent to school. However, as Science and the requirements for citizenship became more a part of education, girls were excluded from learning these topics. Higher education continued to develop, with the 1746 opening of The College of New Jersey (later known as Princeton), and King's College (now Columbia) in 1754. All of these universities were meant exclusively for White men, though some of the colleges experimented by admitting Native Americans. In the public schools, vocational education expanded. Though what was lost by failing to educate the underclasses was incalculable, we can gauge the lost possibilities through such individuals as Benjamin Banneker and Phillis Wheatley. Mr. Banneker, a self-educated free African-American, observed the stars, wrote his own almanac, and was one of the surveyors of what would later become the District of Columbia. Miss Wheatley, an African-born slave educated and freed by her mistress, wrote a remarkable volume of poems published in the year 1773. Most of these had been published in "The Newport Mercury", edited by Benjamin Franklin's brother James. Questions For Review. 1. What were the reasons for the French and Indian War? 2. What was the strategy of General Braddock against the French at Fort Duquesne? What was the strategy of the defending French and Indian forces? 3. Examine the succession of acts imposed upon the American Colonists in the wake of the war, beginning with the Sugar Act. What was the intended purpose of each act? What was its actual effect? 

Background. The British forces might at first glance seem to have every advantage. At the outset of the War they had stocks of cannon and ammunition. The Colonists had single-shot rifles from local forges, guns which took time to load and could easily misfire or explode. When Washington took command of the army in 1775, he learned that there was only enough gunpowder to provide nine rounds of ammunition per man. The British had a large professional army drilled to a pitch like that of Ancient Rome, well-supplied with food, uniforms, and arms. But the American lack of training meant that they did not mass in the European style. Instead they relied on snipers, individuals hidden behind the trees who shot their bullet and then loaded again while their neighbors fired. They had learned this during the French and Indian War. Snipers helped strengthen the American odds. The Beginning of the War (1775 - 1778). Lexington and Concord. The British government commanded General Thomas Gage to enforce the Intolerable Acts and limit rights in Massachusetts. Gage decided to confiscate a stockpile of colonial arms located in Concord. On April 19, 1775, Gage's troops marched to Concord. On the way, at the town of Lexington, Americans who had been warned in advance by Paul Revere and others of the British movements made an attempt to stop the troops. No one knows which side fired the first shot, but it sparked battle on Lexington Green between the British and the Minutemen. Faced against an overwhelmingly superior number of British regular troops in an open field, the Minutemen were quickly routed. Nevertheless, alarms sounded through the countryside. The colonial militias poured in and were able to launch guerrilla attacks on the British while they marched on to Concord. The colonials amassed of troops at Concord. They engaged the British in force there, and they were able to repulse them. They then claimed the contents of the armory. The British retreated to Boston under a constant and withering fire from all sides. Only a reinforcing column with artillery support on the outskirts of Boston prevented the British withdrawal from becoming a total rout. The following day the British woke up to find Boston surrounded by 20,000 armed colonists, occupying the neck of land extending to the peninsula the city stood on. The Battle of Bunker Hill. The action changed from a "battle" to a "siege", where one army bottles up another in a town or a city. (Though in traditional terms, the British were not besieged, since the Royal Navy controlled the harbor and supplies came in by ship.) General Artemas Ward, the head of the Massachusetts militia, had the initial oversight of the siege. He set up headquarters at Cambridge, Massachusetts and positioned his forces at Charlestown Neck, Roxbury, and Dorchester Heights. The 6,000 to 8,000 rebels faced some 4,000 British regulars under General Thomas Gage. Boston and little else was controlled by British troops. General Gage countered the siege on June 17 by attacking the colonists on Breed's Hill and Bunker Hill. Although the British suffered tremendous casualties compared to the colonial losses, the British were eventually able to dislodge the American forces from their entrenched positions. The colonists were forced to retreat when many colonial soldiers ran out of ammunition. Soon after, the area surrounding Boston fell to the British. However, because of the losses they suffered, they were unable to break the siege of the city. Despite the early defeat for the colonists, the battle proved that they had the potential to counter British forces, which were at that time considered the best in the world. The Last Chance For Peace. The Second Continental Congress adopted the Olive Branch Petition, a petition for peace, on July 5, 1775. The Congress affirmed its allegiance to the Crown. It was received in London at the same time as it heard of the Battle For Bunker Hill. The King refused to read the petition or to meet with its ambassadors. Parliament reacted by passing the Prohibitory Act, which banned trade with the colonies. Battle For Boston. Despite the British access to the ships, the town and the army were on short rations. Salt pork was the order of the day, and prices escalated rapidly. While the American forces had some information about what was happening in the city, General Gage had no effective intelligence of rebel activities. On May 25, 1775, 4,500 reinforcements and three new generals arrived in Boston Harbor. The fresh leaders were Major General William Howe and Brigadiers John Burgoyne and Henry Clinton. Gage began planing to break out of the city. On July 3, 1775, George Washington arrived to take charge of the new Continental Army. Forces and supplies came in from as far away as Maryland. Trenches were built at Dorchester Neck, extending toward Boston. Washington reoccupied Bunker Hill and Breeds Hill without opposition. However, these activities had little effect on the British occupation. In the winter of 1775– 1776, Henry Knox and his engineers under order from George Washington used sledges to retrieve sixty tons of heavy artillery that had been captured at Fort Ticonderoga. Knox, who had come up with the idea to use sledges, believed that he would have the artillery there in eighteen days. It took six weeks to bring them across the frozen Connecticut River, and they arrived back at Cambridge on January 24, 1776. Weeks later, in an amazing feat of deception and mobility, Washington moved artillery and several thousand men overnight to take Dorchester Heights overlooking Boston. General John Thomas fortified the area. The British fleet had become a liability, anchored in a shallow harbor with limited maneuverability, and under the American guns on Dorchester Heights. When General Howe saw the cannons, he knew he could not hold the city. He asked that George Washington let them evacuate the city in peace. In return, they would not burn the city to the ground. Washington agreed: he had no choice. He had artillery guns, but did not have the gunpowder. The whole plan had been a masterful bluff. The siege ended when the British set sail for Halifax, Nova Scotia on March 17, 1776. The militia went home, and in April Washington took most of the Continental Army forces to fortify New York City. Ethan Allen and Fort Ticonderoga. The British had considered Fort Ticonderoga a relatively unimportant outpost in a conflict which had up to then been mostly based in Massachusetts. However, a veteran of the French and Indian War, Ethan Allen, had his eye on the fort. Allen had built up a Vermont territorial militia, the Green Mountain Boys, until it was an effective fighting force. Vermont was claimed by the New York colony, but Allen wanted more independence. In April of 1775, Allen was surprised by a visit by Commander Benedict Arnold of the Connecticut Militia. Arnold announced that he had been commissioned to seize the cannons of Fort Ticonderoga. A heated discussion between the two concluded with the agreement that the two militias would combine to attack the fort. This was for the best, for both forces together were small, well short of brigade strength. On May tenth, the combined American forces captured the fort. They seized the arms, including the cannons, which were then hauled by oxen all the way to Boston. Strengthening The Cause. Through the media available in that day, the Revolution promoted the idea of honorable men in revolt against tyranny. Newspapers in North and South published incendiary stories and inspiring engravings. The theater contributed dramatic outcries, including those of Mary Otis Warren. Songs were played and sung to rally flagging spirits. Thomas Paine's 1776. In January of 1776, the Englishman Thomas Paine published the pamphlet "Common Sense". This anti-monarchical publication encouraged American independence, using examples from the Bible and republican virtues to argue that kings were never good for any free state. In late 1776 he began printing his series of pamphlets, "The American Crisis", calling soldiers to mass to the cause of the Revolution. The first of these pamphlets begins with the stirring words, "These are the times that try men's souls." The Declaration of Independence. As military hostilities built up, the Second Continental Congress appointed George Washington as General of the Continental Army. Washington gave up his salary for the position all through the war. (As he was among the richest men in the colonies, he could afford this choice.) In June of 1776, the Second Continental Congress felt it needed a spur for separation from Great Britain. It appointed a Committee of Five to draft a declaration of independence: John Adams, Benjamin Franklin, Robert Livingston, Roger Sherman and Thomas Jefferson. Jefferson became the principal author of this document. Although the British king was no longer principally responsible for his dominion's policy, the Declaration of Independence called him a tyrant. It justified the rights of the rebellion with words the European Enlightenment would have hailed: "We hold these truths to be self-evident, that all men are created equal[.]" The Continental Congress signed the document on July 4, 1776. However, the signatures at this point showed that they wished independence: they could not alone achieve it. Army Bands. One of the attributes of a well-drilled company of soldiers was its military band. The British and Hessian troops drilled to the beat of drums, which carried the rhythm of the march above the noise of musket fire, and provided a way to communicate on the battlefield. "By 1778 soldiers marched at seventy-five 24″ steps per minute in common time and nearly double that (120 steps per minute) when marching in quick time." The music of the fife (a shrill flute) and drum helped build soldier morale. If the Rebellion could not have good supplies, it would at least have high morale. With the "Middle-brook Order", George Washington directed that every officer must provide military music for his troops. This was despite the limited number of instruments. The bands were used to announce the beginning and end of the day, direct troops in battle, and uplift spirits. Popular Music in the Revolution. One of the two major songs of the Revolution was the hymn "Chester", first published in 1770 in "The New England Psalm Singer" and revised in 1778. Its author and composer, William Billings, created a combination of the biblical ("Let tyrants shake their iron rod") and the topical ("Howe and Burgoyne and Clinton too,/ With Prescot and Cornwallis join'd"). Another was the song "Yankee Doodle", adapted from a tune of the Seven Years War. This was originally used by the British to laugh at the provincial manners of the Colonists, but was turned into a theme of the American upstarts. Canada. In September of 1775, the Colonists, led by General Richard Montgomery, invaded Canada. At first the invasion proved successful, with Montgomery capturing Fort St. Jean and the city of Montreal. On December 30 he made the decision to launch an attack onto the British held city of Quebec. It proved disastrous, and Montgomery was killed in battle. This was the last major action in Canada, although Benidict Arnold and a number of other generals did attack the coasts or Canada, or launch raids across the border. The Turning Point of the War. Despite the numerous defeats they faced in the early years of the war, the colonists were able to turn the tide around with several major victories. New York and New Jersey. In July, 1776, General William Howe and thirty-thousand British troops arrived at Staten Island in New York. The large army attacked and defeated General George Washington's American forces in the Battle of Long Island. After nearly having his entire army captured, Washington led a skilled withdrawal out of New York. Eventually the Continental Army was forced to set up camp in Pennsylvania. Howe could have ended the war by pursuing Washington's forces. But Howe was very cautious and took almost no risks. He feared losing too many men so far from home. Britain hired German mercenaries (Hessians) to guard the British fort at Trenton. Howe took advantage of these replacements and decided to wait until spring to attack the Continental Army again. Washington also took advantage of the situation, though from a different perspective. He figured that the Hessians would be weakest on Christmas night, after heavy feasting and drinking. On the night of December 25, 1776, Washington led his troops 9 miles, and across the Delaware River to ambush the Hessians. Crossing the river was difficult. A hail and sleet storm had broken out early in the crossing, winds were strong and the river was full of ice floes. The crossing took 3 hours longer than expected, but Washington decided to continue the attack anyway. As Washington predicted, the mercenaries were completely caught off guard and had little time to respond. Within just a over an hour, on the morning of December 26, the Continental Army had won the Battle of Trenton. The Americans had just 4 wounded and 0 killed against 25 Hessians Killed, 90 wounded and 920 captured. The victory increased the troops' morale and eventually led to re-enlistments. Some historians even speculate Trenton saved the revolution. On January 2, the British came to re-take Trenton, and did so with heavy casualties. Washington once again led a clever withdrawal, and advanced on Princeton. At the Battle of Princeton, the Continental Army attacked the rear-guard of the British Army, and forced them to retreat from New Jersey. During the war, the New World was being devastated by the 1775–1782 North American smallpox epidemic. Having survived the diseased in his youth, and having been warned about the effect the disease may have on the Army by Benjamin Franklin, George Washington wrote he had more dread of the disease crippling the continental army then the British Troops. On February 5, George Washington ordered the first mass inoculation of troops, following large disruptions caused by smallpox outbreaks. The policy was unpopular among soldiers, but stopped the mass infections from continuing. The Battle of Saratoga. In the summer of 1777, British General John Burgoyne and General Howe agreed to attack the colonial Army from two sides and defeat it. Howe marched north, winning the Battles of Brandywine and Germantown and eventually capturing Philadelphia. But Burgoyne was not so fortunate. Delayed by natural traps set up by the Continental Army, his troops slowly marched from Canada to Albany. By September of the year, his forces reached Saratoga, where an enormous American Army attacked the troops. In October, General Burgoyne surrendered all his forces to the Americans. General Howe resigned his post, thwarted despite his victories in Pennsylvania. The Battle of Saratoga proved to be the major turning point in the war. It persuaded France that America had to overthrow Great Britain, and French aid now was introduced to the colonists. The battle was also the last time the British would advance North. By the summer of 1778, following the Battle of Monmouth in New Jersey, all fighting would take place in the South. Defeat of the Iroquois. The Iroquois Confederacy in its zenith had been the equal of the European Powers. But since the French and Indian war it had been in decline. The Tribes of the Confederacy disagreed on who to support in the Revolution. The Onedia and Tuscaroras supported the Americans, while the Mohawk, Onondaga, Cayuga, and the Seneca supported the British. The Confederacy managed to stay together until 1777, when following the Battle of Saratoga, the 4 Tribes supporting the British began to attack American settlements across New York and Pennslyvenia. A back and forth battle followed. The Iroquois would attack American Forts and Towns, then the Americans would burn Iroquois villages. In 1779 George Washington sent General Sullivan to destroy the Iroquois Nation. After defeating the Iroquois at the Battle of Newtown, Sullivan's army then carried out a scorched earth campaign, methodically destroying at least forty Iroquois villages. The devastation created great hardships for the thousands of Iroquois refugees outside Fort Niagara that winter, and many starved or froze to death. The survivors fled to British regions in Canada and the Niagara Falls and Buffalo areas. Thus ended the 700-year history of the Iroquois Confederacy. Conclusion of the War (1778 - 1781). After the loss at Saratoga, the French, traditional rivals of the British, offered their aid in the Revolution. The United States allied itself with France in 1778. Spain and the Dutch Republic also joined the American side, both lending money to the United States and going to war with Britain. Valley Forge. Following the capture of Philadelphia by the British, Washington took his followers to Valley Forge on December 17th, 1777, a defensible nearby area, and built camp for 12,000 soldiers and 400 civilians, then the fourth largest settlement in the colonies. Following the introduction of the Prussian protege of Frederick the Great, Baron Von Steuben to Congress by way of Benjamin Franklin, he was directed to Valley Forge, where he arrived on February 23, 1778. On the Seas. War broke out on the seas as well. Americans granted commissions to "privateers" to attack and destroy all British ships, whether they were military or not. One of the most famous privateers, John Paul Jones, scored several victories at sea for the Americans, even attacking the shores of Britain itself. The War Heads South. An attempted treachery was defeated when its architect, British Major John Andre, was captured in September of 1780. Benedict Arnold, one of the heroes of Fort Ticonderoga, had been placed in charge of Fort Clinton, New York (now called West Point). In response to a bribe, Arnold neglected maintenance of the fortification, and was then preparing to turn the fort over to the British. After he had learned of Andre's arrest he fled to join the British army. Britain turned its attention from the North to the South, where more loyalists lived. They were at first very successful, defeating the Americans at Waxhaws, Charleston, and Camden. Lord Cornwallis, commander of the British forces in the south, was faced with the challenge of chasing down the Americans. Nathanael Greene had split his army into two, leaving one under the control of Daniel Morgan. Morgan drew Banastre Tarleton, who was commanding one half of the British Army, to Cowpens where they were they decisively defeated the British. The other half of the British Army, still under control of Cornwallis, defeated the Americans at the Battle of Guilford Court House. However, it was a bloody victory for Cornwallis and he was forced to withdraw to Yorktown Virginia to regroup. After hearing that the British were in Yorktown, and there was a French Fleet arriving, Washington took the Continental Army, along with French Troops, to Yorktown and surrounded the British. By mid September the town was under siege. Cornwallis was assured by British Commander-in-Chief, Henry Clinton, who was in New York, that he would be relieved shortly. However, the British relief force was defeated by the French fleet. The British continued to hold off for a few more days, but the allied army moved in closer and closer to Yorktown, and their cannons destroyed many of the British defenses. On October 19, 1781, Cornwallis surrendered his entire army, over 7,000 men. Scattered fighting continued, but back in Britain, the British were crushed by this defeat. Parliament voted to cease all offensive operations in "the colonies." Washington took his army to Newburgh, New York, where he stopped a mutiny in the Army. At the conclusion of the war in 1783 large numbers of loyalists and their families relocated to the home country of England and in large part to Canada as well as to other British Colonies. They submitted claims for lost property and lands in America. Many of the claims were not accepted by the English government for lack of evidence of the losses or significantly reduced. The property and lands were acquired by the American communities and then resold to the highest bidders. Due to the climatic effects of a 1782 eruption of an Icelandic volcano, the loyalists also experienced one of the coldest Canadian winters on record which contributed to poor crops in 1783-1784. Starvation, disease and hardship were rampant and many resolved to return to the United States despite the threats of retribution rather than subsist on their meager produce. Treaty of Paris (1783). The British lost hope of crushing the rebellion after Yorktown. They decided to negotiate peace with The United States, France, and Spain. The Treaty of Paris was signed on September 3rd, 1783. In it, the United States was recognized as an independent nation, with boundaries stretching from the Canadian border in the North, to the Mississippi River to the West, and the northern border of Florida in the South. Britain was forced to return Florida to Spain, but could still hold Canada. Congress was told to advise the states to restore property lost or stolen from the Loyalists. (However, many Loyalists had fled during the Revolution, and many of them did not return to claim their property.) Religion &amp; the Revolution. Catholics in the Revolution. The complex situation of Catholicism in Great Britain had results in its Colonies. At the time of the American revolution, Catholics formed approximately 1.6% of the total American population of the original 13 colonies. If Catholics were seen as potential enemies of the British state, Irish Catholics, subject to British rule, were doubly-damned. In Ireland they had been subject to British domination. In America Catholics were still forbidden from settling in some of the colonies. Although the head of their faith dwelt in Rome, they were under the official representation of the Catholic Bishop of the London diocese, one James Talbot. When War began, Bishop Talbot declared his faithfulness to the British Crown. (If he had done otherwise, Catholics in England would have been in trouble. Anti-Catholic sentiment still ran high.) He forbade any Colonial priest to serve Communion. This made practice of the faith impossible. This created sympathy for the Colonial rebels. The Continental Army's alliance with the French increased sympathy for the faith. When the French fleet arrived in Newport, Rhode Island, the colony repealed the Act of 1664 and allowed citizenship to Catholics. (This anticipated the provision of the Constitutional Bill of Rights which would strike anti-Catholic laws from the books.) After the war, the Pope created an American Bishop, John Carroll -- a descendant of the same Carrolls who had helped found Maryland -- and an American Diocese communicating directly with Rome. From Anglicanism to Episcopalianism. On the one hand the colonial Church of England was an organ of the British government and a collaborator with it. Its clergy swore allegiance to the King. Several colonial governments paid monies to the local Anglican Church. Although other faiths were allowed in those states, the Anglican was considered the Official (Established) Church, putting pressure on other denominations. Still, several Revolutionaries, including Thomas Jefferson, rented their pews in a Church of England building. (Jefferson's own faith was Low Church, and he disagreed with the miracles in Christianity.) Significant meetings of the rebellion were held in Church of England buildings. But after the war, the Church needed to find a new role. Some of the Loyalist clergy went north to Canada. Others were allowed to remain after swearing an oath to the new government. The formerly Established Church was no more: even before the creation of the Constitution, with its separation of Church and State, Americans did not want to pay any extra fees. The Book of Common Prayer, the form of worship in that Church, was pragmatically revised for the new Episcopal Church so that people prayed for "Civil Rulers," instead of the King. But many Church buildings were closed, and there was now room for other denominations to flourish in Virginia and other states. The Early Government of the New United States. [Copied from Wikipedia] The Articles of Confederation, formally the Articles of Confederation and Perpetual Union, was an agreement among the 13 founding states that established the United States of America as a confederation of sovereign states and served as its first constitution. Its drafting by the Continental Congress began in mid-1776, and an approved version was sent to the states for ratification in late 1777. The formal ratification by all 13 states was completed in early 1781. Even when not yet ratified, the Articles provided domestic and international legitimacy for the Continental Congress to direct the American Revolutionary War, conduct diplomacy with Europe and deal with territorial issues and Native American relations. Nevertheless, the weakness of the government created by the Articles became a matter of concern for key nationalists. [Whom?] Questions for Review. 1. Who were these authors/composers, and what were they known for? (Mary Otis Warren, Thomas Jefferson, Thomas Paine, William Billings.) 2. How do the battles of Lexington and Concord show the early strengths and weaknesses of the American fighters? 3. Examine a copy of the Declaration of Independence in relation to this and the previous chapters. How does its rhetoric (choice of words) address the concerns of the American rebellion? How does it deviate from actual events to make a point? 

Ideas and Questions of the Time. The overriding question throughout the decade preceding the Civil War was, “Should slavery be allowed in the new territories of the United States?” Before 1848, the question had been hypothetical; however, with the new lands acquired during the Mexican War, it was time for America to make a firm decision regarding the expansion of slavery. The central ideas dominating the debate were: The Wilmot Proviso. On August 8, 1846, Representative David Wilmot, a Pennsylvania Democrat, presented a proposal expressing that “slavery nor involuntary servitude shall ever exist in any part of any territory obtained from Mexico.” The Wilmot Proviso was never accepted as law, but it at long last put the issue forth on the political table. The Calhoun Resolutions. John C. Calhoun, the South Carolina statesman, responded with the Calhoun Resolutions, which said that Congress had no right to stop any citizen with slaves in their possession from taking those slaves into one of the territories. If they did so, the Fifth Amendment, which states that no person can be “deprived of life, liberty, or property, without due process of law,” would be violated. While this was not made formal legislation either, this belief became the standard in most of the south. Popular Sovereignty. A third option, which appealed to many moderates, most prominently Stephen A. Douglas of Illinois, was the idea of popular sovereignty. This was the idea of letting the settlers of a territory themselves decide whether slavery was to be allowed in it, by voting on state constitutions and other such measures. The primary merit of this initiative was that it took the debate out of Congress, which quickly grew tired of the issue, and put it into the hands of people it truly affected. There was also an unspoken understanding that most of the territories would end up being free, as most settlers that were already in those areas did not bring their slaves with them. Compromise of 1850. America looked to the Senate for an answer to the question of slavery within the territories. Henry Clay, nicknamed the "Great Compromiser," constructed a compromise: California was admitted as a free state, but all other territories in the Mexican Cession were allowed to choose between becoming a free territory or a slave territory. Also, as part of the Compromise, the slave trade was banned in the District of Columbia, and a Fugitive Slave Act was passed to allow the capture of fugitive slaves. The Fugitive Slave Act was a very controversial measure. Previously, many in the North felt that slavery merely occurred in the South and that they had nothing to do with it. But under the Fugitive Slave Act, Northerners were required to help return runaway slaves. Thus, the Northerners felt that they were being dragged into aiding the institution of slavery. Several Northern states passed laws prohibiting their officials from aiding the enforcement of the Act. While the admission of California as a free state gave the free states the majority in Congress, the pro-slavery measures in the Fugitive Slave Act made the Compromise seem more favorable to the South. Uncle Tom’s Cabin. Harriet Beecher Stowe’s Uncle Tom’s Cabin, published in 1852, is often called “the book that started the Civil War.” The melodramatic story of the evil overseer Simon Legree and his slaves Eliza and Uncle Tom painted an accurate picture of the horrors of slavery, and gave rise to much abolitionist feeling in the North. However, the effects were not easily visible from the start: because the country was growing tired of the sectional bickering over slavery, it took a while for the story to becoming embedded in the American imagination. Nat Turner. Nat, commonly called Nat Turner, (October 2, 1800 – November 11, 1831) was an American slave whose slave rebellion in Southampton County, Virginia, was the most remarkable instance of black resistance to enslavement in the antebellum southern United States. His methodical slaughter of white civilians during the uprising makes his legacy controversial, but he is still considered by many to be a heroic figure of black resistance to oppression. At birth he was not given a surname, but was recorded solely by his given name, Nat. In accordance with a common practice, he was often called by the surname of his owner, Samuel Turner. Election of 1852. In one of the less spectacular elections in American history, Senator Franklin Pierce of the Democratic party defeated General Winfield Scott of the Whig party. The Whigs tried to rely on Scott’s heroics as a general during the Mexican war to get him elected, a strategy that proved unsuccessful. Pierce, of New Hampshire, ended up being largely an ineffective president, trying and failing to please both the North and the South. The Kansas-Nebraska Act and its Effects. Throughout this time, plans were underway for a transcontinental railroad. A question arose as to what Eastern city should be the main terminus. Senator Stephen Douglas of Illinois hoped to advance his own state’s interests by making Chicago the railroad hub. To do this, he suggested a piece of legislation known as the “Kansas-Nebraska Act,” requiring recognition of two new territories, Kansas and Nebraska, west of Missouri and Iowa, respectively. These territories would both help his railroad and solve the overdue issue of the territories in the remainder of the Louisiana Purchase. But to get the Kansas-Nebraska Act passed, he would have to get the support of Southerners, who wanted a railroad along a more southern route. For this reason, Douglas included in the Act the provision of popular sovereignty in the territories. This blatantly violated the Missouri Compromise of 1821, which stated that slavery would be prohibited above the 36º30’ line. Douglas therefore opened himself up to the verbal barrage of protests from the North, who denounced the cancellation of the Missouri Compromise as unfair. Yet the Act passed, to the indignation of many Northerners, with the support of President Pierce. The North. Many in the North figured that if the Missouri Compromise was not an unbreakable law, neither was the Fugitive Slave Act, leading to many demonstrations against it. Boston witnessed the most remarkable of these, leading to many New Englanders turning against Pierce for his support of the Kansas-Nebraska Act. Political Parties. The Whig party essentially buckled under the pressure of the Kansas-Nebraska Act, with the North condemning it and the South supporting it. Whigs from the North joined some Democrats and Free Soilers that united under the general principle of the Wilmot Proviso, eventually calling themselves the Republican Party and offering its first presidential candidate, John C. Fremont in 1856. Rachel v. Walker. Rachel v. Walker was a lawsuit involving a slave who, in 1834, sued for her freedom from John Walker in the Supreme Court of Missouri, and won. This result was cited in 1856 in the famous Dred Scott v. Sandford case before the Supreme Court of the United States.[1] Dred Scott. The question of the constitutionality of Congressional Compromises was decided by the Supreme Court in 1856. In "Scott v. Sanford", the Court ruled against a slave, Dred Scott, who had sued to become free. The Court ruled 7-2 that Scott remained a slave, and there were nine written opinions. The Chief Justice of the United States, Roger Taney, decided that blacks were so inferior that they could not be citizens of the United States, and that, consequently, they could not sue for his freedom (a state issue)in diversity in federal court, and therefore the court lacked jurisdiction. Nevetheless (the biggest "nevertheless" in American history) in a supererogatory effort to settle the question of slavery once and for all, the Marylander Taney ruled that the Missouri Compromise (which had banned the expansion of slavery into the territories north of Missouri) among other laws, was unconstitutional because it restricted the Constitutional right to own property. Many felt that Taney had committed a legal error in his decision. First, Taney had ruled that Scott had no right to sue. The case should have ended there. Taney had ruled on the constitutionality of the Missouri Compromise, which had, under Taney's own ruling that Scott had no right to sue, no bearing upon the case. Thus, the outrage against the Dred Scott decision was increased even more. John Brown’s Raid. John Brown, an extreme abolitionist last seen engineering the Pottawatomie Massacre in Kansas, came to the federal arsenal at Harper’s Ferry, Virginia for his last fight. He planned to take over the arsenal, give weapons to the slaves that would support him, and make a center of black power in the Appalachian Mountains that would support slave uprisings in the south. The raid did not go quite as planned. Brown did take over the arsenal and took a couple of hostages, but ended up being assaulted by Virginia militia and U.S. Marines under the command of Col. Robert E. Lee of the US 2nd Cavalry. He was tried, convicted, and hanged for treason to the State of Virginia. However, his Raid left a profound impact. John Brown became a martyr for the abolitionist cause during the Civil War. In the South, his actions gave cause to rumors of Northern conspiracy supporting slave insurrections, engendering further suspicion of outsiders in the South. A later Northern marching song sang “John Brown’s body lies a-mouldering in the grave, but his soul is marching on.” Lincoln.  &lt;br&gt; " Lincoln campaign poster" In 1860, four major candidates ran for President. The Whigs, adopting the name "Constitutional Union", nominated Tennessean Senator John Bell. The Northern Democrats nominated Senator Stephen Douglas of Illinois and the Southern Democrats nominated the Vice President John Breckenridge of Kentucky. The more united Republican party nominated Abraham Lincoln, who spoke out against expansion of slavery. Though he assumed that, under the constitution, Congress could not outlaw slavery in the South, he assured all that he would work to admit only free states to the US. Due to divisions between the parties, Lincoln won the election by carrying every Northern State. Douglas won Missouri, Bell the Upper South, and Breckenridge the Deep South. The South was outraged. The North had a far larger population than the South, and thus had more electoral votes. The South had been out voted. 

This textbook is based on the College Entrance Examination Board test in Advanced Placement United States History. The test is a standard on the subject, covering what most students in the United States study in high school and college, so we treat it as the best reference. The text was reorganized and edited in November 2008 to be closer to the content and organization the college board requires. The content was carefully chosen for significance and interest. We welcome reader feedback and suggestions for improvement. Enjoy! The AP Course Description can be found here. 

GCSE Science/Electricity Now we have learned Ohm's Law we can start applying it to some circuits. The test circuits below can be used to investigate how the voltage across a component varies with the current flowing through the component. Either circuit can be used. The one on the left is easiest to set up, the one on the right is easiest to use once it is set up, but does require you to use the rheostat as a potential divider. Note that the voltmeter is in parallel with the component and the ammeter is in series with it.This is necessary so that the ammeter and voltmeter do not interfere with the circuit in anyway. An ideal meter does not change either the current or the voltage. The wire should be "resistance wire", such as nichrome. Ordinary copper wire would short circuit the power pack and blow the pack's fuse. Set up either circuit and then by turning the voltage setting on the power pack and the slider on the rheostat, adjust the current to read 0.1 A. Take a reading of the voltage. Repeat for 0.2 A, 0.3 A up to 1.0 A. Repeat the readings but this time go from 1.0 A to 0.9 A and so on down to 0.1 A. Set your results out in a table like this. Now plot the average voltage against the current on graph paper. You should find the points all fall approximately on a straight line. The slope of the line =Voltage/Current. This is the resistance. Below is a typical voltage/current graph for a wire. Note that all the points fall approximately on a straight line. The slope of the line is 1.1 V/A. Therefore the resistance of this wire is 1.1 Ω All "Ohmic" components have a constant resistance like this. However, "non-Ohmic" components (that is, components that do not obey Ohm's law) do not have such a resistance. A bulb, for example, has a resistance that increases as the current flowing through it goes up, because the filament is heating up. A current-voltage graph for a bulb would be a curve. Q1) Plot a current/voltage graph for the following set of data for a 12V filament bulb. Answers | «Current, Voltage, Resistance | Parallel and series circuits» 

Complex numbers are the extension of the real numbers, i.e., the number line, into a number plane. They allow us to turn the rules of plane geometry into arithmetic. Complex numbers have fundamental importance in describing the laws of the universe at the subatomic level, including the propagation of light and quantum mechanics. They also have practical uses in many fields, including signal processing and electrical engineering. Introduction. Currently, we are able to solve many different kinds of equations for formula_1 , such as formula_2 , or formula_3 , or formula_4 . In each of these cases the solution for formula_1 is a real number: respectively 5, 4/3, and -10. However, there is no real number "x" that satisfies the equation formula_6 , since the square of any real number is nonnegative. Conceptually it would be nice to have some kind of number to be the solution of formula_6 . This "number" would not be a real number, however, and we refer to such a number as an "imaginary" number. Now, is this really the reason? Well - Definitely not! This mistake occurs in many teaching books from the attempts to "solve problems by force", as we could explain psychologically. This has nothing to do with reality, and gives the false feeling that mathematicians are "Cranks" full of to much spare time on their hands with nothing to do. The reason, be surprised, has to do with a problem called Cubic functions. We then extend the real number system to accommodate this special number. It turns out that there will be two imaginary solutions of the equation formula_6 . One of them will be called formula_9 and, following the normal rules for arithmetic, the other solution is formula_10 . We may be inclined to say that formula_11 . That would, however, be incorrect solely because in words this says that "the square root of -1 is formula_9" , but there is no basis for preferring formula_9 over formula_10 (or vice versa) as the square root of -1. Rather, the two square roots have equal standing. We say that all numbers of the form a +bformula_9 , where formula_16 and formula_17 are any real numbers, is the set of "complex" numbers, and we denote this set formula_18 . The real numbers formula_19 may be considered to be the subset of complex numbers formula_20 for which b = 0. Complex numbers can be added, subtracted, multiplied, and divided (except by 0). We will explore some of the properties of these numbers later. There are in fact two commonly used definitions of complex numbers, but they are immediately seen to be logically equivalent. For a negative root like formula_21 , we split the number into two parts such that one part is formula_22 like formula_23 which leads to formula_24 Definition 1. A complex number is an expression of the form "x + yi", in which "x" and "y" are real numbers and "i" is a new number, called the imaginary unit, for which expressions the normal rules of calculation apply together with the extra rule: "i2=-1". Definition 2. A complex number is a pair of real numbers "(x,y)", satisfying the properties: In both cases a complex number consists of two real numbers "x" and "y". The real number "x" is called the real part and the real number "y" the imaginary part of the complex number. From the properties we deduce that complex numbers of the form (x,0) behave just like the real numbers, so we identify (1,0) with 1 and hence (x,0) with x. Furthermore we see that: It is common use to write "i" instead of (0,1), so: Any complex number "(x,y)" may now be written as "x + yi". Some examples. A "complex number" is a number that is in the form formula_30 , where "a" and "b" are real numbers. We say that "a" is the "real" part of z and write formula_31, and that "b" is the "imaginary" part of z, and write formula_32 A number of the form "b"i is sometimes called a "pure imaginary" number, as it has no real part. The pure imaginary numbers are also complex numbers, because "b"i = 0 + "b"i. In the same way, all real numbers are also complex numbers, because "a" = "a" + 0i. So the set of complex numbers includes real numbers, pure imaginary numbers, and the sums of reals and pure imaginaries. Here are some examples Notice that the number 2 is a complex number and a real number. This fact is clearer if we write 2 = 2 + 0i. Any complex number may be written in three main forms, which we will explore later. The form "x" + "y"i is known as the "Cartesian" form. Complex numbers and matrices. Complex numbers can be identified with a certain set of "matrices". If we think of the 2×2 identity matrix as the number 1, and we think of formula_9 which we introduced above as the matrix then the complex number formula_35 then has the form Properties. Complex numbers obey most of the properties of real numbers. Take two complex numbers, formula_37 and formula_38. Addition. How can we add these two complex numbers? We don't even have to think about formula_9 as being "special" in any way, just treat it as any other symbol and proceed by the standard rules of algebra, grouping along the way. We obtain: If one uses the matrix analogy above, regular matrix addition works to add complex numbers in the same way. Verify for yourself that this is true. Subtraction. Subtraction proceeds just as before. Multiplication. By the normal rules, taking into account that formula_42, we find: If one uses the matrix analogy above, regular matrix multiplication works to multiply complex numbers in the same way. Verify for yourself that this is true. Conjugates. The "conjugate" of a complex number formula_44, written formula_45, is the same number with the sign of the imaginary part changed: the conjugate of "a" + "b"i = "a" - "b"i (and vice versa). Let us examine what happens when we have a complex number "z" = "a" + "b"i, what is the product of "z" and its conjugate? Notice the imaginary parts cancel out, so the product is a "real number". This will aid us greatly in the division of a complex number, as we will see. Notice also that this is the "sum" of two squares, analogous to the difference of two squares. Matrix transposition behaves as conjugation if one uses the matrix analogy. Division. How do we compute the quotient of two complex numbers? It is not difficult. Let the quotient be: then cross multiplication gives: hence and So we have to solve two linear equations. Solution: and Note that this is complete nonsense unless formula_53. In fact, it is easy to see that the pair of linear equations have a solution exactly when this is true, i.e., when c and d are not both zero. Or in other words, when formula_54 as a complex number. Thus we can divide by a complex number "only" when it is non-zero. Luckily there is a little trick to speed up this computation. We multiply the denominator with a well chosen number as to make it real. We "realize the denominator" by multiplying with the conjugate of the denominator; a complex number times its conjugate is a real number: hence: Note that in the multiplication and division of complex numbers, we usually work out the whole problem instead of just memorizing the equation of the answer. The reader will note that we use here the familiar trick from algebra of multiplying by the number 1 in a particularly convenient form: formula_57. (We leave it to the reader to verify that any non-zero complex number divided by itself is in fact equal to 1.) Exponents and Roots. Because formula_58, real numbers can be raised to imaginary numbers. Imaginary and complex numbers cannot be raised to imaginary or complex numbers, as imaginary and complex numbers have no natural logarithms. Example: formula_59 Because formula_60, imaginary numbers can be root degrees. Problem set. Given the above rules, answer the following questions. Note: Use sqrt(x) for formula_61 &lt;quiz display=simple points="2/2"&gt; (7 + 2i) + (11 - 6i) = { 18_19 } + { -4_19 } i (8 - 3i) - (6i) = { 8_19 } + { -9_19 } i (9 + 4i)(3 - 16i) = { 91_19 } + { -132_19 } i 3i formula_62 9i = { -27_19 } + { 0_19 } i formula_63{ 1/5|0.2_19 } + { 2/5|0.4_19 } i formula_64{ [11sqrt(3) - 12]/19 (i)|[11sqrt(3)-12]/19 (i) _19 } + { [44 + 3sqrt(3)]/19 (i)|[44+3sqrt(3)]/19 (i) _19 } i formula_65{ x/(x^2 + y^2) (i)|x/(x^2+y^2) (i) _19 } + { -y(x^2 + y^2) (i)|-y(x^2+y^2) (i) _19 } i &lt;/quiz&gt; The Argand Plane. We can represent complex numbers "geometrically" as well. Every complex number can be represented in the form z=x+iy (so x=Re(z) and y=Im(z)). Then we can represent z in the xy-plane by the point (x,y). Notice that this is a one-to-one relationship: for each complex number, we have one corresponding point in the plane, and for each point in the plane there corresponds one complex number. When we use the xy-plane in this way to represent complex numbers, we call the plane the "Argand plane". We will refer to the "y" axis as the imaginary axis, and the "x" axis as the real axis. Notice that a purely imaginary number is represented in the Argand plane by a point on the imaginary axis. A purely real number is represented by a point on the real axis. Here is an example of the Argand plane. There are two complex numbers drawn in the plane; viz., 1 + i, -2 - i. Their sum is plotted on the graph, -1. The red and blue lines demonstrate how geometrically, a parallelogram can be constructed, and their apex forms their sum. Modulus and argument. On this diagram we can see the number 3 + 4i. The red line is the distance away from the origin (the number 0 + 0i). The gray line represents the distances away from the respective axes. We can see that the red line makes an angle θ from the real axis. It is clear that almost all complex numbers have this distance away from the origin and that almost all complex numbers make an angle away from the real axis. We give these two qualities special names; the distance away from the origin is known as the "modulus" of the complex number, and the angle θ is known as the "argument" of the complex number. We write the modulus of a complex number z by |z|, and the argument of the complex number as arg z. We can calculate the modulus and argument by basic trigonometry. Calculating the modulus. In the above example, we have the number 3 + 4i. We can form a triangle in the Argand plane with base 3 and height 4. By Pythagoras we can find the length of the hypotenuse by formula_66. And thus, the length of the hypotenuse is thus the modulus of the complex number, and it is 5 for 3+4i. Generalization. If "z" = "x" + "y"i, |"z"| is clearly formula_67. Equivalently, formula_68. Calculating the argument. We have the same triangle as we had in calculating the modulus. Remember from trigonometry that tan θ is the ratio of the height over the base. So, for 3 + 4i, we have tan θ = 4/3, and thus θ = arctan 4/3 = 0.9... With complex numbers, we always take two things: Note that arg 0 is undefined. Generalization. If "z" = "x" + "y"i, arg "z" is clearly arctan (y/x), or, equivalently, arg "z" = arctan (Im("z")/Re("z")). The polar form. We are now able to calculate the modulus and argument of a complex number, where these two numbers are able to uniquely describe every number in the Argand plane. Using these two characteristics of complex numbers, we are now able to formulate a new way of writing these numbers. Note that in the above diagram, we obtain a triangle that describes the complex number 3 + 4i. Clearly, we can do this for all complex numbers in the Argand plane (except for 0). To simplify our work, let us look at numbers in the circle of unit length equidistant from 0. From trigonometry, we can parameterize all the points on a circle in the Cartesian plane by (cos θ, sin θ). In complex number notation we can say that all numbers on this unit circle are in the form cos θ+i sin θ. This works well on the unit circle, but how does this generalize to describing "all" numbers on the plane? We simply make the circle larger or smaller to encompass the number; this is done by multiplying by the modulus. So then, we obtain the polar form "r"(cos θ + i sin θ) = "z", where "r" is the modulus. Euler's formula. A very significant result in the area of complex numbers is Euler's formula. It basically asserts that This statement can be verified through a rearrangement of the of the cosine and sine functions. Note that conjugate complex numbers have an opposing argument. 2e2i and 2e-2i are conjugate pairs. Proofs. Using Taylor series. Here is a proof of Euler's formula using expansions as well as basic facts about the powers of "i": The functions "e""x", cos("x") and sin("x") (assuming "x" is real) can be written as: and for complex "z" we "define" each of these function by the above series, replacing "x" with "iz". This is possible because the radius of convergence of each series is infinite. We then find that The rearrangement of terms is justified because each series is absolutely convergent. Taking "z" = "x" to be a real number, gives the original identity as Euler discovered it.formula_77 Using calculus. Define the complex number "z" such that Differentiating "z" with respect to "x": Using the fact that "i"2 = -1: Separating variables and integrating both sides: where C is the constant of integration. To finish the proof we have to argue that it is zero. This is easily done by substituting x = 0. But z is just equal to: thus So now we just exponentiate Corollaries. A number of significant results follow as corollaries to Euler's result. de Moivre's theorem. De Moivre's theorem is useful in calculating powers of complex numbers. It states that This follows clearly (from the laws of exponents) if we rewrite the theorem in the form Equivalent trigonometric forms. From the cosine/sine form of the complex number, we can rewrite the cosine and sine function in terms of exponentials. Relation of e, formula_91, i, 1, and 0. By substituting π into the formula, we obtain the following result: formula_92 The actual mathematical relevance of this equation is actually very little. It is known more for its relation of the many branches of mathematics: e comes from calculus, π from geometry, i comes from algebra, 1 is the multiplicative identity, and 0 is the additive identity. It is also known for its simple mathematical aesthetic. Forming trigonometric identities. The aforementioned equivalent trigonometric forms, combined with the Binomial theorem, allow us to create some trigonometric identities that would be difficult to form in any other way. These identities can be used to simplify integral problems. Cosine/sine powers. How can we simplify, say, (cos "x")5? Let us look at a simple example for motivation. First, rewrite as 1/25 (eix+e-ix)5 By the Binomial theorem, Replacing a with ei"x"/2 and b with e-i"x"/2, we obtain Procedure. We can summarize from the above example the general procedure: To simplify an expression in the form: the procedure is: Cosine/sine multiples. We can also form identities in the form: Let's look at another example to see how it's done. Example. Let's expand sin(3"x"). Recall de Moivre's theorem stating We will use this fact to expand out the left side. For ease of manipulation, it may be easier to let "c" = cos("x"), and "s"=sin("x"). Then use the Binomial theorem again to expand out: Collect real and imaginary parts Now, this is of course equal to cos(3"x")+i sin (3"x"). So, substituting back cos("x") for "c" and similarly for sin("x"), we can equate real and imaginary parts, and we get and we get for free. NB: In the cosine expansion, one could write sin("x")2 as 1-cos("x")2 to obtain a formula consisting completely of cosines. (analogously one could write the sine expansion using only sines) External links. This is incomplete and a draft, additional information is to be added. 

Lesson 20, as a bit of a reward is a little translation exercise from the Gospel of Saint Luke. Exercise 1 Vocabulary coming soon, at the moment consult your dictionary Respondens Simon dixit: "Aestimo quia is, cui plus donavit". At ille dixit ei: "Recte iudicasti". Et conversus ad mulierem, dixit Simoni: "Vides hanc mulierem? Intravi in domum tuam: aquam pedibus meis non dedisti; haec autem lacrimis rigavit pedes meos et capillis suis tersit. Osculum mihi non dedisti haec autem, ex quo intravi non cessavit osculari pedes meos. Oleo caput meum non unxisti; haec autem unguento unxit pedes meos. Propter quod dico tibi: Remissa sunt peccata eius multa, quoniam dilexit multrum: cui autem minus dimittitur, minus diligit." Dixit autem ad illam: "Remissa sunt peccata tua". Et coeperunt, qui simul accumbebant, dicere intra se: "quis est hic, qui etiam peccata dimittit?". Dixit autem ad mulierem: Fides tua te salvam fecit; vade in pace!". Et factum est deinceps, et ipse iter faciebat per civitatem et castellum praedicans et evangelizans regnum Dei, et Duodecim cum illo, et mulieres aliquae, quae erant curatae ab spiritibus malignis et infirmitatibus, Maria, quae vocatur Magdalene, de qua daemonia septem exierant, et Ioanna uxor Chuza procuratoris Herodis, et Sussanna et aliae multae, quae ministrabant eis de facultatibus suis. 

Congratulations! You have completed the Introductory Latin course! To finish off, you can look at an exciting section from Caesar's Civil Wars Book II stanza 10–12. Exercise 1 Skill: Translation from Latin source "Vocabulary", coming soon Ubi ex ea turri, quae circum essent opera, tueri se posse confisi sunt, musculum pedes LX longum ex turrim murumque perducerent, facere instituerunt cuius musculi haec erat forma. Duae primum trabes in solo aeque longae distantes inter se pedes IIII collocantur, inque eis columellae pedum in altitudinem V defiguntur. Has inter se capreolis molli fastigio coniungunt, ubi tigna, quae musculi tegendi causa ponant, collocentur. Eo super tigna bipedalia iniciunt eaque laminis clabisque religant. Ad extremum musculi tectum trabesque extremas quadratas regulas IIII patentes digitos difigunt, quae lateres, qui super musculo struantur, contineant. Ita fastigato atque ordinatim structo, ut trabes erant in capreolis collocatae. lateribus lutoque musculus, ut ab igni, qui ex muro iaceretur, tutus esset, contegitus. Super lateres coria inducuuntur, ne canalibus aqua immissa lateres diluere posset. Coria autem, ne rursus igni ac laidibus corrumpantur, centonibus conteguntur. Hoc opus omne tectum vineis ad ipsam turrim perficiunt subitoque inoopinantibus hostibus machinatione navali, phalangis subiectis, ad turrim hostium admovent, ut aedificio iungatur. Quo malo perterriti subito oppiindani saxa quam maxima possunt vectibus promovent praecipitataque muro in musculum devolvunt. Ictum firmitas materiae sustinet, et quicquid incidit fastigio musculi elabitur. Id ubi vident, mutant consilium: cupas taeda ac pice refertas incendunt easque de muro musculum devolvunt. Involutae labuntur, delapsae ab lateribus longuriis furcisque abopere removentur. Interim sub musculo milites vectibus infima saxa turris hostium, quibus fundamenta continebantur convellunt. Musculus ex turri latericia a nostris telis tormentisque defenditur; hostes ex muro ac turribus submoventur: non datur libera muro defendendi facultas. Compluribus iam lapidibus ex ea, quae suberat, turri subductis repentina ruina pars eius turrisconcidit, pars reliqua consequens procumbebat: cum hostes urbis direptione perterriti inermes cum infulis se porta foras universi proripiunt, ad legatos atque exercitum supplices manus tendunt. 

For your pleasure, I present to you a selection of poems from Catullus. Tips for translation: soon.  Exercise 1  Skill: Reading Latin poetry  I  Cui dono lepidum novum libellum  arido modo pumice expolitum?  Corneli, tibi : namque tu solebas  meas esse aliquid putare nugas,  iam tum cum ausus es unus Italorum  omne aevum tribus explicare chartis  doctis, Juppiter, et laboriosis.  quare habe tibi quicquid hoc libelli,  qualecumque, quod, o patrona virgo,  plus uno maneat perenne seclo.  Exercise 2  Skill: Reading Latin poetry  V  Vivamus, mea Lesbia, atque amemus  rumoresque senum severiorum  omnes unius aestimemus assis.  soles occidere et redire possunt.  nobis, cum semel occidit brevis lux  nox est perpetua una dormienda.  da mi basia mille, deinde centum,  dein mille altera, dein secunda centum  deinde usque altera mille, deinde centum;  dein, cum milia multa fecerimus  conturbabimus illa, ne sciamus,  aut nequis malus invidere possit,  cum tantum sciat esse basiorum.  Exercise 3  Skill: Reading Latin poetry  XVI  Pedicabo ego vos et irrumabo,  Aureli pathice et cinaede Furi,  qui me ex versiculis meis putastis,  quod sunt molliculi, parum pudicum,  nam castum esse decet pium poetam  ipsum, versiculos nihil necessest.  qui tum denique habent salem ac leporem.  si sunt molliculi ac parum pudici  et quod pruriat incitare possunt,  non dico pueris, sed hi s pilosis,  qui duros nequeunt moveere lumbos.  vos, quod milia multa basiorum  legistis, male me marem putatis?  pedicabo ego vos et irrumabo.  Exercise 4  Skill: Reading Latin poetry  XCVII  Non (ita me di ament) quicquam referre putavi,  utrumne os an culum olgacerem Aemilio.  nilo mundius hoc, niloque immundior ille,  verum etiam culus mundior et melior,  nam sine dentibus est: os dentis sesquipedalis,  gingivas, vero ploxeni habet veteris,  praeterea rictum qualem diffissus in aestu  meientis mulae cunnus habere solet.  hic futuit multas et se facit esse venustum,  et non pistrino traditur atque asino?  quem siqua attingit, non illam posse putemus  aegroti culum lingere carnificis? 

GCSE Science/Electricity You will have already studied series and parallel circuits before, so should be familiar with what a series and parallel circuit is and their basic properties.But in this module we will be looking at them in a little more detail. We will apply Ohm's law to see how we can work out the resistance of a whole circuit that is made up of a large number of components. Before we begin You might want to try some revision questions. Follow the link then come back here when you are finished.  GCSE Science/Parallel and series circuits revision questions. Resistors in Series Circuits. As we know, 'in a series circuit the current in all parts of the circuit is the same and the current has a one way system. The current depends on the applied voltage and the number of and nature of other components in the circuit. Consider two resistors in a series circuit with a battery. As you might expect the total resistance in this circuit is higher than the resistance of each resistor, because the battery has to "push" the charge through both resistors one after the other. So the total potential difference of the supply is "shared" between the two resistors. Think about what is happening to the current as it flows around the circuit. There are no branches, nowhere for the electric current to escape to, so obviously the same current must flow through both resistors. Let's call this current I. The total resistance for the whole circuit is very simple. It's just R1 plus R2. Formula for resistance of two resistors in series. RTotal = R1 + R2 The total resistance in a series circuit is the sum of the resistances of all the components. Using this formula we can calculate the voltage across each resistor. But it could depend on how many resistors there are. Example: Calculating voltages in a series circuit. Question: Suppose R1 = 1Ω and R2= 4Ω. If the battery supplies 2.5 Volts what is the voltage across each resistor. Answer: First use the total resistance of the circuit to work out how much current is flowing through the circuit. Step 1: Work out the total resistance Step 2: Use Ohms's law to calculate the current in the circuit. Decide which resistance to use. Now we know how much current is flowing through both resistors, we can work out the voltage across each resistor. Step 3: Use Ohm's law to calculate the voltage across each resistor: Resistors in Parallel Circuits. Parallel circuits are a bit more complicated than series circuits, because they contain a branch - the electric current will take more than one route. Look at the diagram. At point X the current splits into 2 paths, and flows through both resistor R1 and resistor R2. It may not split in equal amounts. When several(two or more) components are connected in parallel branches, the voltage (potential difference) across each parallel branch is the same. And this is the same as the voltage across the battery. The current through each component is the same as if it were the only component present. So the total current flowing through the battery is the sum of the currents flowing through each branch. Here is the formula for the currents flowing through a parallel circuit. IMain = I1 + I2 Because the voltage is the same for all branches, the branch with the lowest resistance has the highest current flowing through it. Because there are more paths for the charge to flow along, the total resistance is less than either of the two paths on their own. And therefore (with the same battery) the current is bigger. To find out the resistance of the whole circuit , we can't just add together the resistors as we did in the series circuit, we have to apply Ohm's law to each branch of the circuit. Imagine we replaced the resistors with bulbs. You should now be able to answer the following questions from your previous knowledge. If you answered the above questions correctly you should find this next section easy! Example: Calculating the resistance of several resistors in parallel. It's best to break the process down into a number of simple steps. Some books may give you a formula to use, but you shouldn't use any formula without understanding where it comes from (otherwise you are likely to remember it incorrectly or apply it inappropriately). By using a simple step by step method instead you can get a feel for why it works, and you will be far less likely to make a mistake in the exam. Step1 considering each branch on its own, as if the other branch didn't exist use Ohm's law to work out the current flowing through each branch. Step 2 Add the currents together to find out the total amount of current flowing through the whole circuit. Step 3 Apply Ohm's law again to work out the total resistance of the circuit. Example. Suppose V=2V R1 = 1Ω and R2= 1Ω. Applying Ohm's law to the branch containing R1 gives I1=V/R1 =2/1 =2A Applying Ohm's law to the branch containing R2 gives I2=V/R2 =2/1 =2A Total current = I1+I2 = 4A Applying Ohm's law again, to the whole circuit. RTotal =V/Itotal = 2/4 =0.5Ω Notice that the resistance of the whole circuit is "lower" than the resistance of either branch! Practice questions. Answers | Three useful components» 

Soluciones a los ejercicios.&lt;br&gt; "(Solutions to the exercises)" Pronombres personales.&lt;br&gt; "(Personal pronouns.)" Ejercicio 1.&lt;br&gt; Rellena los espacios en blanco con los pronombres personales correctos en español. "Exercise 1."&lt;br&gt; "Fill the blank spaces with the correct Spanish personal pronous." Ejercicio 2.&lt;br&gt; Rellena los espacios en blanco con los pronombres personales correctos en español. "Exercise 2."&lt;br&gt; "Fill the blank spaces with the correct Spanish personal pronouns." Enlace a los Ejercicios "Link to exercises" Enlace a la lección 1 "Link to Lesson 1." 

Soluciones a los ejercicios. "Solutions to the exercises" "Fill in the blank". Por favor, rellena los espacios en blanco con la forma correcta del verbo: "Please, fill in the blank with the correct form of the verb:" 1. Nosotros aprendemos español. "We learn Spanish." 2. Yo compro un libro. "I buy a book." 3. Carmen y Roberto viajan a Mexico. "Carmen and Roberto travel to Mexico." 4. Ana habla ingles. "Ana speaks English." 5. Tú bebes una cerveza. "You drink a beer." 6. Susana escribe una carta. "Susana writes a letter." 7. Los niños estudian para el examen. "The children study for the exam." 8. Fernando y Lucas cantan una cancion. "Fernando and Lucas sing a song." 9. Tú lees un libro. "You read a book." 10. Vosotros subís las escaleras. "You (plural) climb the stairs." Enlace a los Ejercicios "Link to exercices" Enlace a la lección 3 "Link to the lesson 3." 

Further Discussion. The questions posed in Chapter 1. Introduction are discussed further on this page. Remember, some questions are intended to be thought-provoking and more than one answer may be "correct". Users are encouraged to expand upon the discussions here.  1-1. "Do you think the "scientific method" is something only a scientist would use?" The "scientific method" is something that is taught, so we might question just how natural or intuitive an approach it is. Think about how you discovered answers to interesting or difficult questions as a child. You might have first sought the opinion of an expert (like your dad, right?). That is the first step used by a scientist as well. Scientists learn to seek the expertise of others that have asked similar questions and come up with at least partial answers. This does not require contacting other scientists; more often it is a matter of becoming very familiar with the literature (the written and published record) on the subject of interest. Scientists rarely formulate "hypotheses" without first learning all they can about a subject from what has been published already. In effect, they learn how to put their question at the very forefront of knowledge about a subject. There are so many questions to be answered in science that no one would want to expend a lot of energy asking and answering questions to mysteries already solved. Of course learning is different. Asking questions and seeking answers from books, journals, and other sources is a valid approach to this first step of learning, and one you and any scientist would have in common. Formulating a proper question or a hypothesis is perhaps the most difficult step. Recall that the hypothesis is not really a question but an answer. The question is the wonder about something; the hypothesis is a testable answer to that question. It is not THE answer, just a reasonable one posed in such a way that an experiment can be conducted to ascertain its validity. You may see this step expressed as "making a guess"; but consider that after he or she has completed the step described above—learning everything written about a subject—the scientist is in a position to make a pretty good guess. This step could be difficult for the scientist because he or she must have in mind one or more experiments to conduct to test the hypothesis. This step may be difficult for you for the same reason, and because it probably feels odd to answer a question before investigating it. The hypothesis step is included in the scientific method to promote intellectual honesty and to allow others to better understand a scientist's reasoning when reviewing what was done and how a conclusion was reached. Compare these scenarios: Or In the absence of any further report from this technician, only the second procedure provides any information to the next technician faced with a second, similar bomb. He at least knows what not to do! Note also that the hypothesis is falsifiable. Testing demonstrated whether an intact red wire was or was not needed for detonation (apparently it controlled the undetonated state). Framing the hypothesis thusly: "the red wire serves a purpose" cannot be proven false; what test could you devise to satisfactorily demonstrate it has no purpose? It may in fact have no purpose (not in the example above), but eliminating by testing every possible purpose under the sun is an unsatisfactory approach.  1-3. When you catch a cold a virus has infected your body.  "Why do you think there is reason to question whether the virus is living or not?"  After all, if you took some ricinin (a plant poison), you would get very sick,  but no one would suggest the toxin were alive or that the plant had entered your body. Ricinin is an alkaloid and, along with the toxalbumin ricin (a plant protein), constitute extremely toxic substances found in the seeds of the castor bean ("Ricinus communis") plant. Were you to ingest several raw seeds, you would experience nausea and vomiting, stomachache, bloody diarrhea, headache, cold sweat, sleepiness, disorientation, fever, shortage of breath, seizures, and possibly collapse and death. While there are a number of different ways that a virus could enter your body, the most common would be breathing in virus particles while in the presence of an infected person. The outcome of such exposure could be pretty much the same as that described for ricinin ingestion, depending upon the type of virus and your body's ability to respond appropriately to the "infection". Or, perhaps the infection you "catch" is cholera—the "Vibrio cholerae" bacterium. These bacteria are typically ingested in drinking water contaminated by improper sanitation. Again, you might display many of the same symptoms described, including the really unpleasant part about dying. The point here is to consider which of these problems constitute an "attack" by another living organism and which represent simply a poisoning of your living body by a non-living chemical. You cannot use the resulting symptoms as an indication because...; well, because they are just symptoms: what you "feel" as your body reacts to the chemical or biological attack. In the case of ricin, this protein inhibits protein synthesis within the cells of the body. Organ damage results. "Vibrio cholerae" in the body produce a chemical that results in a loss of fluid and salts across the lining of the gut. The resulting diarrhea allows the bacterium to spread to other people under unsanitary conditions, and can result in death due to dehydration (an inability to retain water). So both the non-living chemical and the living bacterium cause illness by toxicity to our bodily functions. You should recognize one pretty significant difference between an illness caused by a living organisms and an illness caused by a non-living, but toxic, chemical. In the latter case, the severity of our symptoms will be pretty much dependent upon the dose of toxin we ingest. There will be just so much chemical entering the body, and the damage should be proportional to that amount. In the case of the bacterium, something like a million cells need to be ingested to result in an infection, but once established, the living organism is a tiny factory that cranks out toxin and reproduces more identical factories (more cells) that themselves produce toxin. So the dose of toxin we get from the bacterial infection is not obviously limited. Living cells metabolize (break down and produce chemicals for various purposes) and reproduce (increase their numbers). In essence, this is what is meant by an infection: another life form has taken up residence in or on our body, and is utilizing organic substances that are a part of our life processes to carry out its life functions, to further its existence and numbers. In the case of chlorea, the bacterium has some nasty habits that make us very ill; but there are a number of other bacteria ("Escherichi coli", for example) that live in perfect harmony within or digetive tract and help us digest food. So not every "alien" that invades our body necessarily causes an infection.  « Chapter 1 

There are two systems of naming months in use: one based on the zodiac, and one based on numbers. The place structure is "x1 is foomonth in year x2 in calendar x3"; this is not according to jvojva but makes it easy to specify dates as "the day nth of the foomonth of the year mmmm". There are also systems of naming the days of the week based on Oriental elements and numbers. Note that Sunday in Lojban is either day seven or day zero, depending on your perspective. 

&lt; Lojban  N north berti B 0°00'  NbE north by east bersunberberti BDBB 11°15'  NNE north northeast bersunberti BDB 22°30'  NEbN northeast by north bersunberstuna BDBD 33°45'  NE northeast berstuna BD 45°00'  NEbE northeast by east bersunsunberti BDDB 56°15'  ENE east northeast bersunstuna BDD 67°30'  EbN east by north bersunsunstuna BDDD 78°45'  E east stuna D 90°00'  EbS east by south nansunsunstuna NDDD 101°15'  ESE east southeast nansunstuna NDD 112°30'  SEbE southeast by east nansunsunsnanu NDDN 123°45'  SE southeast nanstuna ND 135°00'  SEbS southeast by south nansunynanstuna NDND 146°15'  SSE south southeast nansunsnanu NDN 157°30'  SbE south by east nansunynansnanu NDNN 168°45'  S south snanu N 180°00'  SbW south by west nansicnansnanu NVNN 191°15'  SSW south southwest nansicysnanu NVN 202°30'  SWbS southwest by south nansicnanstici NVNV 213°45'  SW southwest nanstici NV 225°00'  SWbW southwest by west nansicysicysnanu NVVN 236°15'  WSW west southwest nansicystici NVV 247°30'  WbS west by south nansicysicystici NVVV 258°45'  W west stici V 270°00'  WbN west by north bersicysicystici BVVV 281°15'  WNW west northwest bersicystici BVV 292°30'  NWbW northwest by west bersicysicyberti BVVB 303°45'  NW northwest berstici BV 315°00'  NWbN northwest by north bersicyberstici BVBV 326°15'  NNW north northwest bersicyberti BVB 337°30'  NbW north by west bersicyberberti BVBB 348°45' 

Lojban. le ninmu cu te dunda le gerku&lt;br&gt; .uinaicaidai zo'e na prami le ninmu&lt;br&gt; .i .a'ucu'idai le ninmu cu na junri zo'e&lt;br&gt; .i .u'icu'idai le ninmu na cmila&lt;br&gt; .i .iucu'idai le ninmu cu na prami zo'e&lt;br&gt; .i .a'odai le nanmu cu klama le ninmu la kentykis.&lt;br&gt; .i .i'odai le nanmu cu dunda le gerku le ninmu&lt;br&gt; .i .oidai le nanmu cu klama la kentykis. le ninmu&lt;br&gt; .i .uiru'edai le gerku cu prami le ninmu&lt;br&gt; .i .a'uru'edai le ninmu cu junri le gerku&lt;br&gt; .i .u'idai le ninmu cu cmila&lt;br&gt; .i .iucaidai le ninmu cu prami le gerku&lt;br&gt; English. The Woman Receives a Dog&lt;br&gt; Sadness! The woman is not loved.&lt;br&gt; Disinterest! The woman doesn't have a gravity for anything.&lt;br&gt; No amusement! The woman doesn't laugh.&lt;br&gt; No love! The woman doesn't love.&lt;br&gt; Hope! The man goes to the woman from Kentucky.&lt;br&gt; Appreciation! The man gives the dog to the woman.&lt;br&gt; Complaint! The man goes to Kentucky from the woman.&lt;br&gt; A little happiness! The dog loves the woman.&lt;br&gt; A little interest! The woman has a gravity for the dog.&lt;br&gt; Amusement! The woman laughs.&lt;br&gt; Lots of love! The woman loves the dog. 

The Articles of Confederation. "(The following text is taken from Wikipedia)" The Articles of Confederation and Perpetual Union, also the Articles of Confederation, was the governing constitution of the alliance of thirteen independent and sovereign states styled "United States of America." The Article's ratification (proposed in 1777) was completed in 1782, legally uniting the states by compact into the "United States of America" as a union with a confederation government. Under the Articles (and the succeeding Constitution) the states retained sovereignty over all governmental functions not specifically deputed to the confederation. The final draft of the Articles was written in the summer of 1777 and adopted by the Second Continental Congress on November 15, 1777 in York, Pennsylvania after a year of debate. In practice the final draft of the Articles served as the "de facto" system of government used by the Congress ("the United States in Congress assembled") until it became "de jure" by final ratification on March 1, 1781; at which point Congress became the Congress of the Confederation. The "Articles" set the rules for operations of the "United States" confederation. The confederation was capable of making war, negotiating diplomatic agreements, and resolving issues regarding the western territories; it could mint coins and borrow inside and outside the United States. An important element of the Articles was that Article XIII stipulated that "their provisions shall be inviolably observed by every state" and "the Union shall be perpetual." This article was put to the test in the American Civil War. The Articles were created by the chosen representatives of the states in the Second Continental Congress out of a perceived need to have "a plan of confederacy for securing the freedom, sovereignty, and independence of the United States." Although serving a crucial role in the attainment of nationhood for the thirteen states, a group of reformers, known as "federalists", felt that the Articles lacked the necessary provisions for a sufficiently effective government. Fundamentally, a federation was sought to replace the confederation. The key criticism by those who favored a more powerful central state (i.e. the federalists) was that the government (i.e. the Congress of the Confederation) lacked taxing authority; it had to request funds from the states. Another criticism of the Articles was that they did not strike the right balance between large and small states in the legislative decision making process. Due to its "one-state, one-vote" plank, the larger states were expected to contribute more but had only one vote. The Articles were replaced by the United States Constitution on June 21, 1788. Background. The political push for the colonies to increase cooperation began in the French and Indian Wars in the mid 1750s. The opening of the American Revolutionary War in 1775 induced the various states to cooperate in seceding from the British Empire. The Second Continental Congress starting 1775 acted as the confederation organ that ran the war. Congress presented the Articles for enactment by the states in 1777, while prosecuting the American Revolutionary war against the Kingdom of Great Britain. Ratification. Congress began to move for ratification of the Articles in 1777: "The articles can always be candidly reviewed under a sense of the difficulty of combining in one general system the various sentiments and interests of a continent divided into so many sovereign and independent communities, under a conviction of the absolute necessity of uniting all our councils and all our strength, to maintain and defend our common liberties..." The document could not become officially effective until it was ratified by all of the thirteen colonies. The first state to ratify was Virginia on December 16, 1777. The process dragged on for several years, stalled by the refusal of some states to rescind their claims to land in the West. Maryland was the last holdout; it refused to go along until Virginia and New York agreed to cede their claims in the Ohio River valley. A little over three years passed before Maryland's ratification on March 1, 1781. Article summaries. Even though the Articles of Confederation and the Constitution were established by many of the same people, the two documents were very different. The original five-paged Articles contained thirteen articles, a conclusion, and a signatory section. The following list contains short summaries of each of the thirteen articles. Still at war with the Kingdom of Great Britain, the colonists were reluctant to establish another powerful national government. Jealously guarding their new independence, members of the Continental Congress created a loosely-structured unicameral legislature that protected the liberty of the individual states. While calling on Congress to regulate military and monetary affairs, for example, the Articles of Confederation provided no mechanism to force the states to comply with requests for troops or revenue. At times, this left the military in a precarious position, as George Washington wrote in a 1781 letter to the governor of Massachusetts, John Hancock. The end of the war. The Treaty of Paris (1783), which ended hostilities with Great Britain, languished in Congress for months because state representatives failed to attend sessions of the national legislature. Yet Congress had no power to enforce attendance. Writing to George Clinton in September 1783, George Washington complained: Function. The Articles supported the Congressional direction of the Continental Army, and allowed the 13 states to present a unified front when dealing with the European powers. As a tool to build a centralized war-making government, they were largely a failure, but since guerrilla warfare was correct strategy in a war against the British Empire's army, this "failure" succeeded in winning independence. Under the articles, Congress could make decisions, but had no power to enforce them. There was a requirement for unanimous approval before any modifications could be made to the Articles. Because the majority of lawmaking rested with the states, the central government was also kept limited. Congress was denied the power of taxation: it could only request money from the states. The states did not generally comply with the requests in full, leaving the confederation chronically short of funds. Congress was also denied the power to regulate commerce, and as a result, the states fought over trade as well. The states and the national congress had both incurred debts during the war, and how to pay the debts became a major issue. Some states paid off their debts; however, the centralizers favored federal assumption of states' debts. Nevertheless, the Congress of the Confederation did take two actions with lasting impact. The Land Ordinance of 1785 established the general land survey and ownership provisions used throughout later American expansion. The Northwest Ordinance of 1787 noted the agreement of the original states to give up western land claims and cleared the way for the entry of new states. Once the war was won, the Continental Army was largely disbanded. A very small national force was maintained to man frontier forts and protect against Indian attacks. Meanwhile, each of the states had an army (or militia), and 11 of them had navies. The wartime promises of bounties and land grants to be paid for service were not being met. In 1783, Washington defused the Newburgh conspiracy, but riots by unpaid Pennsylvania veterans forced the Congress to leave Philadelphia temporarily. Signatures. The Second Continental Congress approved the Articles for distribution to the states on November 15 1777. A copy was made for each state and one was kept by the Congress. The copies sent to the states for ratification were unsigned, and a cover letter had only the signatures of Henry Laurens and Charles Thomson, who were the president and secretary to the Congress. But, the "Articles" at that time were unsigned, and the date was blank. Congress began the signing process by examining their copy of the "Articles" on June 27 1778. They ordered a final copy prepared (the one in the National Archives), and that delegates should inform the secretary of their authority for ratification. On July 9, 1778, the prepared copy was ready. They dated it, and began to sign. They also requested each of the remaining states to notify its delegation when ratification was completed. On that date, delegates present from New Hampshire, Massachusetts, Rhode Island, Connecticut, New York, Pennsylvania, Virginia and South Carolina signed the Articles to indicate that their states had ratified. New Jersey, Delaware and Maryland could not, since their states had not ratified. North Carolina and Georgia also didn't sign that day, since their delegations were absent. After the first signing, some delegates signed at the next meeting they attended. For example, John Wentworth of New Hampshire added his name on August 8. John Penn was the first of North Carolina's delegates to arrive (on July 10), and the delegation signed the "Articles" on July 21 1778. The other states had to wait until they ratified the "Articles" and notified their Congressional delegation. Georgia signed on July 24, New Jersey on November 26, and Delaware on February 12 1779. Maryland refused to ratify the "Articles" until every state had ceded its western land claims. On February 2, 1781, the much-awaited decision was taken by the Maryland General Assembly in Annapolis. As the last piece of business during the afternoon Session, "among engrossed Bills" was "signed and sealed by Governor Thomas Sim Lee in the Senate Chamber, in the presence of the members of both Houses… an Act to empower the delegates of this state in Congress to subscribe and ratify the articles of confederation" and perpetual union among the states. The Senate then adjourned "to the first Monday in August next." The decision of Maryland to ratify the Articles was reported to the Continental Congress on February 12. The formal signing of the "Articles" by the Maryland delegates took place in Philadelphia at noon time on March 1, 1781 and was celebrated in the afternoon. With these events, the Articles entered into force and the United States came into being as a united, sovereign and national state. Congress had debated the "Articles" for over a year and a half, and the ratification process had taken nearly three and a half years. Many participants in the original debates were no longer delegates, and some of the signers had only recently arrived. The "Articles of Confederation and Perpetual Union" were signed by a group of men who were never present in the Congress at the same time. The signers and the states they represented were: Presidents of the Congress. The following list is of those who led the Congress of the Confederation under the "Articles of Confederation" as the presidents of the United States in Congress Assembled. Under the Articles, the president was the presiding officer of Congress, chaired the Cabinet (the Committee of the States) when Congress was in recess, and performed other administrative functions. He was not, however, a "chief" executive in the way the successor "President of the United States" is a chief executive, but all of the functions he executed were under the auspices and in service of the Congress. "For a full list of presidents of the Congress Assembled and presidents under the two Continental Congresses before the Articles, see President of the Continental Congress." Revision and replacement. In May 1786, Charles Pinckney of South Carolina proposed that Congress revise the Articles of Confederation. Recommended changes included granting Congress power over foreign and domestic commerce, and providing means for Congress to collect money from state treasuries. Unanimous approval was necessary to make the alterations, however, and Congress failed to reach a consensus. In September, five states assembled in the Annapolis Convention to discuss adjustments that would improve commerce. Under their chairman, Alexander Hamilton, they invited state representatives to convene in Philadelphia to discuss improvements to the federal government. Although the states' representatives to the Constitutional Convention in Philadelphia were only authorized to amend the Articles, the representatives held secret, closed-door sessions and wrote a new constitution. The new Constitution gave much more power to the central government, but characterization of the result is disputed. Historian Forrest McDonald, using the ideas of James Madison from "Federalist 39", describes the change this way: Historian Ralph Ketcham comments on the opinions of Patrick Henry, George Mason, and other antifederalists who were not so eager to give up the local autonomy won by the revolution: According to their own terms for modification (Article XIII), the Articles would still have been in effect until 1790, the year in which the last of the 13 states ratified the new Constitution. The Congress under the Articles continued to sit until November 1788, overseeing the adoption of the new Constitution by the states, and setting elections. Historians have given many reasons for the perceived need to replace the articles in 1787. Jillson and Wilson (1994) point to the financial weakness as well as the norms, rules and institutional structures of the Congress, and the propensity to divide along sectional lines. Rakove (1988) identifies several factors that explain the collapse of the Confederation. The lack of compulsory direct taxation power was objectionable to those wanting a strong centralized state or expecting to benefit from such power. It could not collect customs after the war because tariffs were vetoed by Rhode Island. Rakove concludes that their failure to implement national measures "stemmed not from a heady sense of independence but rather from the enormous difficulties that all the states encountered in collecting taxes, mustering men, and gathering supplies from a war-weary populace." The second group of factors Rakove identified derived from the substantive nature of the problems the Continental Congress confronted after 1783, especially the inability to create a strong foreign policy. Finally, the Confederation's lack of coercive power reduced the likelihood for profit to be made by political means, thus potential rulers were uninspired to seek power. When the war ended in 1783, certain special interests had incentives to create a new "merchant state," much like the British state people had rebelled against. In particular, holders of war scrip and land speculators wanted a central government to pay off scrip at face value and to legalize western land holdings with disputed claims. Also, manufacturers wanted a high tariff as a barrier to foreign goods, but competition among states made this impossible without a central government. Historical importance. The Articles are historically important for two major reasons: "i)" they were the first constitution or governing document for the United States of America and "ii)" they legally established a union of the thirteen founding states; a Perpetual Union. Early on, tensions developed surrounding the Union, not least because of the fact that with the US Constitution the basis of government was changed from that of confederation to federation. Thomas Jefferson and John C. Calhoun were in their time leading proponents of guaranteeing the constitutional rights of states in federal legislation. Over time, a legal view developed that if the union violated the constitutional rights of states they might rightfully seceed. A significant tension in the 19th century surrounded the expansion of slavery (which was generally supported in agricultural Southern states and opposed in industrial Northern states). As the secessionist view gained support in the South, the opposing view in the North was that since the U.S. Constitution declared itself to be "a more perfect union" than the Articles, it too must be perpetual, and also could not be broken without the consent of the other states. This view was promoted by Daniel Webster and Abraham Lincoln. In 1861, these constitutional contracts were cited by President Lincoln against any claims by the seceding states that unilaterally withdrawing from the Union and taking federal property within those states was legal. The Northwest Ordinance. The Congress established the Northwest Territory around the Great Lakes between 1784 and 1787. In 1787, Congress passed the Northwest Ordinance banning slavery in the new Territory. Congressional legislation divided the Territory into "townships" of six square miles each and provided for the sale of land to settlers. The Northwest Territory would eventually become the states of Ohio, Wisconsin, Indiana, Illinois and Michigan. Problems with the Confederation. The Confederation faced several difficulties in its early years. Firstly, Congress became extremely dependent on the states for income. Also, states refused to require its citizens to pay debts to British merchants, straining relations with Great Britain. France prohibited Americans from using the important port of New Orleans, crippling American trade down the Mississippi river. Shays' Rebellion. Due to the post-revolution economic woes, agitated by inflation, many worried of social instability. This was especially true for those in Massachusetts. The legislature's response to the shaky economy was to put emphasis on maintaining a sound currency by paying off the state debt through levying massive taxes. The tax burden hit those with moderate incomes dramatically. The average farmer paid a third of their annual income to these taxes from 1780 to 1786. Those who couldn't pay had their property foreclosed and were thrown into crowded prisons filled with other debtors. But in the summer of 1786, a revolutionary war veteran named Daniel Shays began to organize western communities in Massachusetts to stop foreclosures, with force, by prohibiting the courts from holding their proceedings. Later that fall, Shays marched the newly formed "rebellion" into Springfield to stop the state supreme court from gathering. The state responded with troops sent to suppress the rebellion. After a failed attempt by the rebels to attack the Springfield arsenal, and other small skirmishes, the rebels retreated and then uprising collapsed. Shays retreated to Vermont by 1787. While Daniel Shays was in hiding, the government condemned him to death on the charge of treason. Shays pleaded for his life in a petition which was finally granted by John Hancock on June 17, 1788. With the threat of treason behind him, Shays moved to New York and died September 25, 1825 U.S. presidents before George Washington. Who was the first president of the United States? Ask any school child and they will readily tell you "George Washington." And of course, they would be correct—at least technically. Washington was inaugurated on April 30, 1789, and yet, the United States continually had functioning governments from as early as September 5, 1774, and operated as a confederated nation from as early as July 4, 1776. During that nearly fifteen-year interval, Congress—first the Continental Congress and then later the Confederation Congress—was always moderated by a duly elected president. This officer was known as the "President of the Continental Congress", and later as the "President of the United States, in Congress Assembled". However, the office of President of the Continental Congress had very little relationship to the office of President of the United States beyond the name. The president of the United States is the head of the executive branch of government, while the president of the Continental Congress was merely the chair of a body that most resembled a legislature, although it possessed legislative, executive, and judicial powers. The following brief biographies profile these "forgotten presidents." Peyton Randolph of Virginia (1723–1775) When delegates gathered in Philadelphia for the first Continental Congress, they promptly elected the former King's Attorney of Virginia as the moderator and president of their convocation. He was a propitious choice. He was a legal prodigy—having studied at the Inner Temple in London, served as his native colony's Attorney General, and tutored many of the most able men of the South at William and Mary College—including the young Patrick Henry. His home in Williamsburg was the gathering place for Virginia's legal and political gentry—and it remains a popular attraction in the restored colonial capital. He had served as a delegate in the Virginia House of Burgesses, and had been a commander under William Byrd in the colonial militia. He was a scholar of some renown—having begun a self-guided reading of the classics when he was thirteen. Despite suffering poor health served the Continental Congress as president twice, in 1774 from September 5 to October 21, and then again for a few days in 1775 from May 10 to May 23. He never lived to see independence, yet was numbered among the nation's most revered founders. Henry Middleton (1717–1784) America's second elected president was one of the wealthiest planters in the South, the patriarch of the most powerful families anywhere in the nation. His public spirit was evident from an early age. He was a member of his state's Common House from 1744 to 1747. During the last two years he served as the Speaker. During 1755 he was the King's Commissioner of Indian Affairs. He was a member of the South Carolina Council from 1755 to 1770. His valor in the War with the Cherokees during 1760–1761 earned him wide recognition throughout the colonies—and demonstrated his leadership abilities while under pressure. He was elected as a delegate to the first session of the Continental Congress and when Peyton Randolph was forced to resign the presidency, his peers immediately turned to Middleton to complete the term. He served as the fledgling coalition's president from October 22, 1774, until Randolph was able to resume his duties briefly beginning on May 10, 1775. Afterward, he was a member of the Congressional Council of Safety and helped to establish the young nation's policy toward the encouragement and support of education. In February 1776 he resigned his political involvements in order to prepare his family and lands for what he believed was inevitable war—but he was replaced by his son Arthur who eventually became a signer of both the Declaration of Independence and the Articles of Confederation, served time as an English prisoner of war, and was twice elected Governor of his state. John Hancock (1737–1793) The third president was a patriot, rebel leader, merchant who signed his name into immortality in giant strokes on the Declaration of Independence on July 4, 1776. The boldness of his signature has made it live in American minds as a perfect expression of the strength and freedom—and defiance—of the individual in the face of British tyranny. As President of the Continental Congress during two widely spaced terms—the first from May 24 1775 to October 30 1777 and the second from November 23, 1785, to June 5, 1786—Hancock was the presiding officer when the members approved the Declaration of Independence. Because of his position, it was his official duty to sign the document first—but not necessarily as dramatically as he did. Hancock figured prominently in another historic event—the battle at Lexington: British troops who fought there April 10, 1775, had known Hancock and Samuel Adams were in Lexington and had come there to capture these rebel leaders. And the two would have been captured, if they had not been warned by Paul Revere. As early as 1768, Hancock defied the British by refusing to pay customs charges on the cargo of one of his ships. One of Boston's wealthiest merchants, he was recognized by the citizens, as well as by the British, as a rebel leader—and was elected President of the first Massachusetts Provincial Congress. After he was chosen President of the Continental Congress in 1775, Hancock became known beyond the borders of Massachusetts, and, having served as colonel of the Massachusetts Governor's Guards he hoped to be named commander of the American forces—until John Adams nominated George Washington. In 1778 Hancock was commissioned Major General and took part in an unsuccessful campaign in Rhode Island. But it was as a political leader that his real distinction was earned—as the first Governor of Massachusetts, as President of Congress, and as President of the Massachusetts constitutional ratification convention. He helped win ratification in Massachusetts, gaining enough popular recognition to make him a contender for the newly created Presidency of the United States, but again he saw Washington gain the prize. Like his rival, George Washington, Hancock was a wealthy man who risked much for the cause of independence. He was the wealthiest New Englander supporting the patriotic cause, and, although he lacked the brilliance of John Adams or the capacity to inspire of Samuel Adams, he became one of the foremost leaders of the new nation—perhaps, in part, because he was willing to commit so much at such risk to the cause of freedom. Henry Laurens (1724–1792) The only American president ever to be held as a prisoner of war by a foreign power, Laurens was heralded after he was released as "the father of our country," by no less a personage than George Washington. He was of Huguenot extraction, his ancestors having come to America from France after the revocation of the Edict of Nantes made the Reformed faith illegal. Raised and educated for a life of mercantilism at his home in Charleston, he also had the opportunity to spend more than a year in continental travel. It was while in Europe that he began to write revolutionary pamphlets—gaining him renown as a patriot. He served as vice-president of South Carolina in 1776. He was then elected to the Continental Congress. He succeeded John Hancock as President of the newly independent but war beleaguered United States on November 1, 1777. He served until December 9, 1778, at which time he was appointed Ambassador to the Netherlands. Unfortunately for the cause of the young nation, he was captured by an English warship during his cross-Atlantic voyage and was confined to the Tower of London until the end of the war. After the Battle of Yorktown, the American government regained his freedom in a dramatic prisoner exchange—President Laurens for Lord Cornwallis. Ever the patriot, Laurens continued to serve his nation as one of the three representatives selected to negotiate terms at the Paris Peace Conference in 1782. John Jay (1745–1829) America's first Secretary of State, first Chief Justice of the Supreme Court, one of its first ambassadors, and author of some of the celebrated Federalist Papers, Jay was a Founding Father who, by a quirk of fate, missed signing the Declaration of Independence—at the time of the vote for independence and the signing, he had temporarily left the Continental Congress to serve in New York's revolutionary legislature. Nevertheless, he was chosen by his peers to succeed Henry Laurens as President of the United States—serving a term from December 10, 1778, to September 27, 1779. A conservative New York lawyer who was at first against the idea of independence for the colonies, the aristocratic Jay in 1776 turned into a patriot who was willing to give the next twenty-five years of his life to help establish the new nation. During those years, he won the regard of his peers as a dedicated and accomplished statesman and a man of unwavering principle. In the Continental Congress Jay prepared addresses to the people of Canada and Great Britain. In New York he drafted the State constitution and served as Chief Justice during the war. He was President of the Continental Congress before he undertook the difficult assignment, as ambassador, of trying to gain support and funds from Spain. After helping Franklin, Jefferson, Adams, and Laurens complete peace negotiations in Paris in 1783, Jay returned to become the first Secretary of State, called "Secretary of Foreign Affairs" under the Articles of Confederation. He negotiated valuable commercial treaties with Russia and Morocco, and dealt with the continuing controversy with Britain and Spain over the southern and western boundaries of the United States. He proposed that America and Britain establish a joint commission to arbitrate disputes that remained after the war—a proposal which, though not adopted, influenced the government's use of arbitration and diplomacy in settling later international problems. In this post Jay felt keenly the weakness of the Articles of Confederation and was one of the first to advocate a new governmental compact. He wrote five Federalist Papers supporting the Constitution, and he was a leader in the New York ratification convention. As first Chief Justice of the Supreme Court, Jay made the historic decision that a State could be sued by a citizen from another State, which led to the Eleventh Amendment to the Constitution. On a special mission to London he concluded the "Jay Treaty," which helped avert a renewal of hostilities with Britain but won little popular favor at home—and it is probably for this treaty that this Founding Father is best remembered. Samuel Huntington (1732–1796) An industrious youth who mastered his studies of the law without the advantage of a school, a tutor, or a master—borrowing books and snatching opportunities to read and research between odd jobs—he was one of the greatest self-made men among the Founders. He was also one of the greatest legal minds of the age—all the more remarkable for his lack of advantage as a youth. In 1764, in recognition of his obvious abilities and initiative, he was elected to the General Assembly of Connecticut. The next year he was chosen to serve on the Executive Council. In 1774 he was appointed Associate Judge of the Superior Court and, as a delegate to the Continental Congress, was acknowledged to be a legal scholar of some respect. He served in Congress for five consecutive terms, during the last of which he was elected President. He served in that office from September 28, 1779 until ill health forced him to resign on July 9, 1781. He returned to his home in Connecticut—and as he recuperated, he accepted more Councilor and Bench duties. He again took his seat in Congress in 1783, but left it to become Chief Justice of his state's Superior Court. He was elected Lieutenant Governor in 1785 and Governor in 1786. According to John Jay, he was "the most precisely trained Christian jurists ever to serve his country." Thomas McKean (1734–1817) During his astonishingly varied fifty-year career in public life he held almost every possible position—from deputy county attorney to President of the United States under the Confederation. Besides signing the Declaration of Independence, he contributed significantly to the development and establishment of constitutional government in both his home state of Delaware and the nation. At the Stamp Act Congress he proposed the voting procedure that Congress adopted: that each colony, regardless of size or population, has one vote—the practice adopted by the Continental Congress and the Congress of the Confederation, and the principle of state equality manifest in the composition of the Senate. And as county judge in 1765, he defied the British by ordering his court to work only with documents that did not bear the hated stamps. In June 1776, at the Continental Congress, McKean joined with Caesar Rodney to register Delaware's approval of the Declaration of Independence, over the negative vote of the third Delaware delegate, George Read—permitting it to be "The unanimous declaration of the thirteen United States." And at a special Delaware convention, he drafted the constitution for that State. McKean also helped draft—and signed—the Articles of Confederation. It was during his tenure of service as President—from July 10, 1781 to November 4, 1782—when news arrived from General Washington in October 1781 that the British had surrendered following the Battle of Yorktown. As Chief Justice of the supreme court of Pennsylvania, he contributed to the establishment of the legal system in that State, and, in 1787, he strongly supported the Constitution at the Pennsylvania Ratification Convention, declaring it "the best the world has yet seen." At sixty-five, after over forty years of public service, McKean resigned from his post as Chief Justice. A candidate on the Democratic-Republican ticket in 1799, McKean was elected Governor of Pennsylvania. As Governor, he followed such a strict policy of appointing only fellow Republicans to office that he became the father of the spoils system in America. He served three tempestuous terms as Governor, completing one of the longest continuous careers of public service of any of the Founding Fathers. John Hanson (1715–1783) He was the heir of one of the greatest family traditions in the colonies and became the patriarch of a long line of American patriots—his great grandfather died at Lutzen beside the great King Gustavus Aldophus of Sweden; his grandfather was one of the founders of New Sweden along the Delaware River in Maryland; one of his nephews was the military secretary to George Washington; another was a signer of the Declaration; still another was a signer of the Constitution; yet another was Governor of Maryland during the Revolution; and still another was a member of the first Congress; two sons were killed in action with the Continental Army; a grandson served as a member of Congress under the new Constitution; and another grandson was a Maryland Senator. Thus, even if Hanson had not served as President himself, he would have greatly contributed to the life of the nation through his ancestry and progeny. As a youngster he began a self-guided reading of classics and rather quickly became an acknowledged expert in the juridicalism of Anselm and the practical philosophy of Seneca—both of which were influential in the development of the political philosophy of the great leaders of the Reformation. It was based upon these legal and theological studies that the young planter—his farm, Mulberry Grove was just across the Potomac from Mount Vernon—began to espouse the cause of the patriots. In 1775 he was elected to the Provincial Legislature of Maryland. Then in 1777, he became a member of Congress where he distinguished himself as a brilliant administrator. Thus, he was elected President in 1781. He served in that office from November 5, 1781 until November 3, 1782. He was the first president to serve a full term after the full ratification of the Articles of Confederation—and like so many of the Southern and New England Founders, he was strongly opposed to the Constitution when it was first discussed. He remained a confirmed anti-federalist until his untimely death. Elias Boudinot (1741–1802) He did not sign the Declaration, the Articles, or the Constitution. He did not serve in the Continental Army with distinction. He was not renowned for his legal mind or his political skills. He was instead a man who spent his entire career in foreign diplomacy. He earned the respect of his fellow patriots during the dangerous days following the traitorous action of Benedict Arnold. His deft handling of relations with Canada also earned him great praise. After being elected to the Congress from his home state of New Jersey, he served as the new nation's Secretary for Foreign Affairs—managing the influx of aid from France, Spain, and Holland. The in 1783 he was elected to the Presidency. He served in that office from November 4, 1782 until November 2, 1783. Like so many of the other early presidents, he was a classically trained scholar, of the Reformed faith, and an anti-federalist in political matters. He was the father and grandfather of frontiersmen—and one of his grandchildren and namesakes eventually became a leader of the Cherokee nation in its bid for independence from the sprawling expansion of the United States. Thomas Mifflin (1744–1800) By an ironic sort of providence, Thomas Mifflin served as George Washington's first aide-de-camp at the beginning of the Revolutionary War, and, when the war was over, he was the man, as President of the United States, who accepted Washington's resignation of his commission. In the years between, Mifflin greatly served the cause of freedom—and, apparently, his own cause—while serving as the first Quartermaster General of the Continental Army. He obtained desperately needed supplies for the new army—and was suspected of making excessive profit himself. Although experienced in business and successful in obtaining supplies for the war, Mifflin preferred the front lines, and he distinguished himself in military actions on Long Island and near Philadelphia. Born and reared a Quaker, he was excluded from their meetings for his military activities. A controversial figure, Mifflin lost favor with Washington and was part of the Conway Cabal—a rather notorious plan to replace Washington with General Horatio Gates. And Mifflin narrowly missed court-martial action over his handling of funds by resigning his commission in 1778. In spite of these problems—and of repeated charges that he was a drunkard—Mifflin continued to be elected to positions of responsibility—as President and Governor of Pennsylvania, delegate to the Constitutional Convention, as well as the highest office in the land—where he served from November 3, 1783, to November 29, 1784. Most of Mifflin's significant contributions occurred in his earlier years—in the First and Second Continental Congresses he was firm in his stand for independence and for fighting for it, and he helped obtain both men and supplies for Washington's army in the early critical period. In 1784, as president, he signed the treaty with Great Britain which ended the war. Although a delegate to the Constitutional Convention, he did not make a significant contribution—beyond signing the document. As Governor of Pennsylvania, although he was accused of negligence, he supported improvements of roads, and reformed the State penal and judicial systems. He had gradually become sympathetic to Jefferson's principles regarding states' rights; even so, he directed the Pennsylvania militia to support the Federal tax collectors in the Whiskey Rebellion. In spite of charges of corruption, the affable Mifflin remained a popular figure. A magnetic personality and an effective speaker, he managed to hold a variety of elective offices for almost thirty years of the critical Revolutionary period. Richard Henry Lee (1732–1794) His resolution "that these United Colonies are, and of right ought to be, free and independent States", approved by the Continental Congress July 2, 1776, was the first official act of the United Colonies that set them irrevocably on the road to independence. It was not surprising that it came from Lee's pen—as early as 1768 he proposed the idea of committees of correspondence among the colonies, and in 1774 he proposed that the colonies meet in what became the Continental Congress. From the first, his eye was on independence. A wealthy Virginia planter whose ancestors had been granted extensive lands by King Charles II, Lee disdained the traditional aristocratic role and the aristocratic view. In the House of Burgesses he flatly denounced the practice of slavery. He saw independent America as "an asylum where the unhappy may find solace, and the persecuted repose." In 1764, when news of the proposed Stamp Act reached Virginia, Lee was a member of the committee of the House of Burgesses that drew up an address to the King, an official protest against such a tax. After the tax was established, Lee organized the citizens of his county into the Westmoreland Association, a group pledged to buy no British goods until the Stamp Act was repealed. At the First Continental Congress, Lee persuaded representatives from all the colonies to adopt this non-importation idea, leading to the formation of the Continental Association, which was one of the first steps toward union of the colonies. Lee also proposed to the First Continental Congress that a militia be organized and armed—the year before the first shots were fired at Lexington; but this and other proposals of his were considered too radical—at the time. Three days after Lee introduced his resolution, in June of 1776, he was appointed by Congress to the committee responsible for drafting a declaration of independence, but he was called home when his wife fell ill, and his place was taken by his young protégé, Thomas Jefferson. Thus Lee missed the chance to draft the document—though his influence greatly shaped it and he was able to return in time to sign it. He was elected President—serving from November 30, 1784 to November 22, 1785 when he was succeeded by the second administration of John Hancock. Elected to the Constitutional Convention, Lee refused to attend, but as a member of the Congress of the Confederation, he contributed to another great document, the Northwest Ordinance, which provided for the formation of new States from the Northwest Territory. When the completed Constitution was sent to the States for ratification, Lee opposed it as anti-democratic and anti-Christian. However, as one of Virginia's first Senators, he helped assure passage of the amendments that, he felt, corrected many of the document's gravest faults—the Bill of Rights. He was the great uncle of Robert E. Lee and the scion of a great family tradition. Nathaniel Gorham (1738–1796) Another self-made man, Gorham was one of the many successful Boston merchants who risked all he had for the cause of freedom. He was first elected to the Massachusetts General Court in 1771. His honesty and integrity won his acclaim and was thus among the first delegates chose to serve in the Continental Congress. He remained in public service throughout the war and into the Constitutional period, though his greatest contribution was his call for a stronger central government. But even though he was an avid federalist, he did not believe that the union could—or even should—be maintained peaceably for more than a hundred years. He was convinced that eventually, in order to avoid civil or cultural war, smaller regional interests should pursue an independent course. His support of a new constitution was rooted more in pragmatism than ideology. When John Hancock was unable to complete his second term as President, Gorham was elected to succeed him—serving from June 6, 1786 to February 1, 1787. It was during this time that the Congress actually entertained the idea of asking Prince Henry—the brother of Frederick II of Prussia—and Bonnie Prince Charlie—the leader of the ill-fated Scottish Jacobite Rising and heir of the Stuart royal line—to consider the possibility of establishing a constitutional monarch in America. It was a plan that had much to recommend it but eventually the advocates of republicanism held the day. During the final years of his life, Gorham was concerned with several speculative land deals which nearly cost him his entire fortune. Arthur St. Clair (1734–1818) Born and educated in Edinburgh, Scotland, during the tumultuous days of the final Jacobite Rising and the Tartan Suppression, St. Clair was the only president of the United States born and bred on foreign soil. Though most of his family and friends abandoned their devastated homeland in the years following the Battle of Culloden—after which nearly a third of the land was depopulated through emigration to America—he stayed behind to learn the ways of the hated Hanoverian English in the Royal Navy. His plan was to learn of the enemy's military might in order to fight another day. During the global conflict of the Seven Years' War—generally known as the French and Indian War—he was stationed in the American theater. Afterward, he decided to settle in Pennsylvania where many of his kin had established themselves. His civic-mindedness quickly became apparent: he helped to organize both the New Jersey and the Pennsylvania militias, led the Continental Army's Canadian expedition, and was elected Congress. His long years of training in the enemy camp was finally paying off. He was elected President in 1787—and he served from February 2 of that year until January 21 of the next. Following his term of duty in the highest office in the land, he became the first Governor of the Northwest Territory. Though he briefly supported the idea of creating a constitutional monarchy under the Stuart's Bonnie Prince Charlie, he was a strident Anti-Federalist—believing that the proposed federal constitution would eventually allow for the intrusion of government into virtually every sphere and aspect of life. He even predicted that under the vastly expanded centralized power of the state the taxing powers of bureaucrats and other unelected officials would eventually confiscate as much as a quarter of the income of the citizens—a notion that seemed laughable at the time but that has proven to be ominously modest in light of our current governmental leviathan. St. Clair lived to see the hated English tyrants who destroyed his homeland defeated. But he despaired that his adopted home might actually create similar tyrannies and impose them upon themselves. Cyrus Griffin (1736–1796) Like Peyton Randolph, he was trained in London's Inner Temple to be a lawyer—and thus was counted among his nation's legal elite. Like so many other Virginians, he was an anti-federalist, though he eventually accepted the new Constitution with the promise of the Bill of Rights as a hedge against the establishment of an American monarchy—which still had a good deal of currency. The Articles of Confederation afforded such freedoms that he had become convinced that even with the incumbent loss of liberty, some new form of government would be required. A protégé of George Washington—having worked with him on several speculative land deals in the West—he was a reluctant supporter of the constitutional ratifying process. It was during his term in the office of the presidency—the last before the new national compact went into effect–that ratification was formalized and finalized. He served as the nation's chief executive from January 22, 1788 until George Washington's inauguration on April 30, 1789. 

^Lesson 2^ "Fill in the blank". Por favor, rellena los espacios en blanco con el articulo correcto: "Please, fill in the blank with the correct article:" 1. ____ mesa "(a table)" 2. ____ caballos "(some horses)" 3. ____ gato "(the cat)" 4. ____ ciudad "(the city)" 5. ____ casas "(some houses)" 6. ____ hijos "(the sons)" 7. ____ hijas "(the daughters)" 8. ____ perro "(a dog)" Soluciones a los ejercicios "Solutions to exercices" ^Lesson 2^ 

^Lesson 2^ Translate the following into Spanish. Soluciones a los ejercicios "Solutions to exercices" ^Lesson 2^ 

GCSE Science/Electricity This page looks at three components that are used in electrical circuits. The three components in question are: The diode, the thermistor, and the light dependent resistor. You will see how the components work and how they are used. The Diode. A diode is a device that allows current to flow in one direction but not in the reverse direction. Its circuit symbol consists of a triangle pointing in the direction that current "is" allowed to flow with a line at the point, inside a circle. Diodes are useful for stopping current flowing in the "wrong" direction. Above is a current voltage graph for a typical diode. Note that the diode needs a small voltage (about 0.6V) to work. Below 0.6V in the forward direction very little current flows. At or above 0.6V the resistance drops dramatically and a huge current can flow. In the reverse direction, virtually no current flows no matter how large the voltage is. (Eventually, when the voltage gets high enough, the diode "breaks down" and current is able to flow in the reverse direction. For normal diodes however this "break down voltage" is very high and can be ignored. A diode that is connected in reverse to use the break down effect is called a Zener diode. It is used to provide voltage control in power supplies (called the Zener voltage).) Diodes are useful for turning AC current into DC current (AC/DC converters are called rectifiers). In AC the electricity flows back and forth in cycles. Mains electricity is AC. It cycles 50 times every second. This is fine for something like a lightbulb, where the bulb does not care in which direction the electricity flows. Some devices however can only work off DC. A toy electric train for example would not go with AC. The current would be saying "go forwards", "go backwards" over and over again 50 times a second! With a diode in the circuit the "negative" part of the cycle is stopped. The graph on the left shows AC current. Note that the current cycles positive, negative, positive, and so on. The electricity is flowing to and fro. Putting the current through a diode has the result seen on the right. The positive part of the cycle in unaffected, but the negative is wiped out. The current goes positive, zero, positive. The Thermistor. A thermistor is a just a resistor whose resistance depends on temperature. Most thermistors have lower resistance when the temperature is high. They are used in circuits which are temperature dependent. For example, fire alarms, the temperature gets too high; the alarm sounds. The LDR. A light dependent resistor is a resistor whose resistance depends on the amount of light shining on it. Its resistance gets less the more light shines on it. It can be used in circuits that are light dependent. For example, switching circuits for street lights. The ambient light gets too low; the lamp turns on. It is also known as a photoresistor from photo, the Latin for light. Questions. Safety in Mains circuits» 

GCSE Science/Electricity Mains electricity is potentially dangerous. There are, however, safety features included in plugs. This module looks at how to wire a plug correctly. The three pin plug. Nowadays most appliances are sold with moulded plugs already fitted. Nevertheless, it is still important to understand the correct wiring of a plug because enough of the old plugs still exist. It is also the case when you bring in equipments overseas. British Standard compliant adaptors are not always available for such non-UK plugs. You are very likely to need to change a plug at some time in your life. In the UK mains electricity is 230 V. (In Hong Kong, it is 220 V.) If you were to touch a live wire a current would flow through your body to the ground. This current may be enough to kill you. The cable from the appliance usually consist of three wires, an earth and two other wires, live and neutral, which carry the current to and from the power station (live is from the power station and neutral is back to the power station). The wires are made of copper surrounded by an insulation casing. The casing is made of plastic and is coloured: The three wires are covered by an outer sheath made of plastic. Q1) Use your knowledge of insulators and conductors to explain The plug has the following features: Q2) Why are the pins made of brass and why is the case plastic? The purpose of the parts of a plug. The live and neutral wires. The live and neutral wires carry the current around the circuit. Mains current is A.C. (alternating current); this means that it is going backward and forwards in cycles (clockwise and anticlockwise around the circuit). The frequency of the cycle is 50 hertz (50 times per second). This cycling of current is achieved by varying the voltage on the line wire from about +325V to – you to the earth. This is where the earth wire is included, for your safety. The earth wire connects the case of the appliance out down the flex to the earth pin on its plug. This connection goes into the socket, then inside the wiring of your house down to the earth through the earthing system (not necessarily plumbing). If the live wire were to touch the case a huge current would flow through the earth wire. This would probably blow the fuse and break the circuit (see next section). However even if the fuse doesn't blow the current would still prefer to flow through a wire with low resistance than a human body with relatively higher resistance. Thus the earth wire helps protect you if you touch the case of an appliance that is "live". The earth pin on a plug is longer than the live and neutral pins. This ensures that the earth pin always connects with the socket first. All sockets have shutters which prevent access to the live contacts when there is no plug in the socket. On some sockets these shutters are operated by the earth pin pushing the shutter mechanism down to uncover the line and neutral socket contacts, on other sockets there is a mechanism which opens the shutters when the two live pins are inserted simultaneously, and on other sockets when all three pins are inserted simultaneously. The fuse. A fuse is simply a very thin wire. The wire has quite a low melting point. As current flows through the wire it heats up. If too large a current flows, it melts, breaking the circuit. Fuses are used to protect the flexible lead between the plug and the appliance. If too large a current flows through a lead it may overheat or catch fire. Fuses are unlikely to act quickly enough to prevent human electrocution – their main purpose is to prevent fires due to large currents. Fuses are rated according to how much current they can carry before melting. In plugs fuses are usually 3A (red), 5A (black), or 13A (brown). The correct fuse is the one that matches the current rating of the lead. All plug fuses must comply to British Standard BS1362. The rating and "BS1362" should be explicitly marked on such fuses. Q4) A table lamp usually carries a current of 0.5A. What fuse should be put in the plug: 3A, 5A, or 13A? Q5) An iron usually carries a current of 5.2A. What fuse should be put in the plug: 3A, 5A, or 13A? Q6) A kettle is protected by an earth wire and a 13A fuse. The line wire comes loose and touches the side of the kettle. The fuse blows. Explain why. Q7) Explain why the fuse is always located on the line wire and not the neutral wire? Q8) Describe and Explain what happens in the following scenarios: Answers | «Three useful components | Power» 

"I'm experimenting here with the practicality of having laboratory excercises as part of the Botany Study Guide" Comments invited. Chapter 3. Plant Structure Laboratory. Leaves. &lt;br&gt; Examine this 'marsh purslane' plant and determine the following about its leaves: 3-1. "Leaf arrangement is": &lt;br&gt; 3-2. "Which of the following describes the leaves": 3-3. "Only one of the following statements is true": « Chapter 3 "Answers to Laboratory Questions:" 

Introduction. 'Number theory' is a large encompassing subject in its own right. Here we will examine the key concepts of number theory. Unlike real analysis and calculus which deals with the dense set of real numbers, number theory examines mathematics in discrete sets, such as N or Z. If you are unsure about sets, you may wish to revisit ../Set theory/. Number Theory, the study of the integers, is one of the oldest and richest branches of mathematics. Its basic concepts are those of divisibility, prime numbers, and integer solutions to equations -- all very simple to understand, but immediately giving rise to some of the best known theorems and biggest unsolved problems in mathematics. The Theory of Numbers is also a very interdisciplinary subject. Ideas from combinatorics (the study of counting), algebra, and complex analysis all find their way in, and eventually become essential for understanding parts of number theory. Indeed, the greatest open problem in all mathematics, the Riemann Hypothesis, is deeply tied into Complex Analysis. But never fear, just start right into "Elementary Number Theory", one of the warmest invitations to pure mathematics, and one of the most surprising areas of applied mathematics! Divisibility. Note that in R, Q, and C, we can "divide" freely, except by zero. This property is often known as "closure" -- the quotient of two rationals is again a rational, etc.. However, if we move to performing mathematics purely in a set such as Z, we come into difficulty. This is because, in the integers, the result of a division of two integers might not be another integer. For example, we can of course divide 6 by 2 to get 3, but we "cannot" divide 6 by 5, because the fraction 6/5 is not in the set of integers. However we can introduce a new relation where division is defined. We call this relation "divisibility", and if formula_1 is an integer, we say: Formally, if there exists an integer formula_10 such that formula_11 then we say that formula_2 divides formula_3 and write formula_14. If formula_2 does not divide formula_3 then we write formula_17: Proposition. The following are basic consequences of this definition. Let a, b, and c be integers: Quotient and divisor theorem. For any integer "n" and any "k" &gt; 0, there is a unique "q" and "r" such that: Here n is known as dividend. We call "q" the "quotient", "r" the "remainder", and "k" the "divisor". It is probably easier to recognize this as division by the algebraic re-arrangement: Modular arithmetic. What can we say about the numbers that divide another? Pick the number 8 for example. What is the remainder on dividing 1 by 8? Using the division theorem above We have a notation for the remainders, and can write the above equations as We can also write These notations are all short for So "x" ≡ 1 (mod 8), for example, is the same as saying Observe that the remainder here, in comparing it with the division algorithm is 1. "x" ≡ 1 (mod 8) asks what numbers have the remainder 1 on division by 8? Clearly the solutions are "x"=8×0+1, 8×1+1... = 1, 9, ... Often the entire set of remainders on dividing by "n" - which we say "modulo n" - are interesting to look at. We write this set Zn. Note that this set is finite. The remainder on dividing 9 by 8 is 1 - the same as dividing 1 by 8. So in a sense 9 is really "the same as" 1. In fact, the relation "≡" is an equivalence relation. We call this relation "congruence". Note that the equivalence classes defined by congruence are precisely the elements of Zn. We can find some number "a" modulo "n" (or we say "a" congruent to "n") by finding its decomposition using the division algorithm. Addition, subtraction, and multiplication work in Zn - for example 3 + 6 (mod 8) = 9 (mod 8) = 1 (mod 8). The numbers do look strange but they follow many normal properties such as commutativity and associativity. If we have a number greater than "n" we often reduce it modulo "n" first - again using the division algorithm. For example if we want to find 11+3 mod 8, its often easier to calculate 3 + 3 (mod 8) rather than reducing 14 mod 8. A trick that's often used is that, say, if we have 6 + 7 (mod 8) we can use "negative" numbers instead so the problem becomes -2 + -1 = -3 = 5 (mod 8). We often use the second notation when we want to look at equations involving numbers modulo some "n". For example, we may want to find a number "x" such that We can find solutions by trial substitution (going through all the numbers 0 through 7), but what if the moduli are very large? We will look at a more systematic solution later. Note: we often say that we are working in Zn and use equals signs throughout. Familiarize yourself with the three ways of writing modular equations and expressions. The Consistency of Modular Arithmetic. Let formula_18 denote an arbitrary base. Given an arbitrary integer formula_19, the sequence of integers formula_20 are all congruent to each other modulo formula_21: formula_22 In modular arithmetic, two integers formula_19 and formula_24 that are congruent modulo formula_21 (formula_26) both "represent" the same quantity from formula_27. It should be possible to substitute an arbitrary integer formula_19 in place of integer formula_24 provided that formula_26. This means that: Number Bases. Converting between various number bases is one of the most tedious processes in mathematics. The numbers that are generally used in transactions are all in base-10. This means that there are 10 digits that are used to describe a number. These ten digits are {0,1,2,3,4,5,6,7,8,9}. Similarly, base-4 has 4 digits {0,1,2,3} and base-2 has two digits {0,1}. Base two is sometimes referred to as Binary. There are also bases greater than 10. For these bases, it is customary to use letters to represent digits greater than 10. An example is Base-16 (Hexadecimal). The digits used in this base are {0,1,2,3,4,5,6,7,8,9,A,B,C,D,E,F}. In order to convert between number bases, it is critical that one knows how to divide numbers and find remainders. To convert from decimal to another base one must simply start dividing by the value of the other base, then dividing the result of the first division and overlooking the remainder, and so on until the base is larger than the result (so the result of the division would be a zero). Then the number in the desired base is the remainders read from end to start. The following shows how to convert a number (105) which is in base-10 into base-2. Answer : 1101001 After finishing this process, the remainders are taken and placed in a row (from bottom to top) after the final quotient (1101001, in this example) is shown as the base-2 equivalent of the number 105. To sum up the process, simply take the original number in base 10, and divide that number repeatedly, keeping track of the remainders, until the quotient becomes less than the numerical value of the base. This works when converting any number from base-10 to any base. If there are any letters in the base digits, then use the letters to replace any remainder greater than 9. For example, writing 11(of base-10) in base 14. Answer: B As 11 is a single remainder, it is written as a single digit. Following the pattern {10=A, 11=B, 12=C...35=Z}, write it as B. If you were to write "11" as the answer, it would be wrong, as "11" Base-14 is equal to 15 in base-10! In order to convert from a number in any base back to base ten, the following process should be used: Take the number 3210 (in base-10). In the units place (100), there is a 0. In the tens place (101), there is a 1. In the hundreds place (102), there is a 2. In the thousands place (103), there is a 3. The formula to find the value of the above number is: 3×103 + 2×102 + 1×101 + 0×100 = 3000 + 200 + 10 + 0 = 3210. The process is similar when converting from any base to base-10. For example, take the number 3452 (in base-6). In the units place (60), there is a 2. In the sixths place (61) there is a 5. In the thirty-sixths place (62), there is a 4. In the 216th place (63), there is a 3. The formula to find the value of the above number (in base-10) is: 3×63 + 4×62 + 5×61 + 2×60 = 648 + 144 + 30 + 2 = 824. The value of 3452 (base-6) is 824 in base-10. A more efficient algorithm is to add the left-most digit and multiply by the base, and repeat with the next digit and so on. ((3 * 6 + 4) * 6 + 5) * 6 + 2 = 824 The processes to convert between number bases may seem difficult at first, but become easy if one practices often. Prime numbers. Prime numbers are the building blocks of the integers. A prime number is a positive integer greater than one that has only two divisors: 1, and the number itself. For example, 17 is prime because the only positive integers that divide evenly into it are 1 and 17. The number 6 is not a prime since more than two divsors 1, 2, 3, 6 divide 6. Also, note that 1 is not a prime since 1 has only one divisor. Some prime numbers. The prime numbers as a sequence begin Euclid's Proof that There are Infinitely Many Primes. The Greek mathematician Euclid gave the following elegant proof that there are an infinite number of primes. It relies on the fact that all non-prime numbers --- composites --- have a unique factorization into primes. Euclid's proof works by contradiction: we will assume that there are a finite number of primes, and show that we can derive a logically contradictory fact. So here we go. First, we assume that that there are a finite number of primes: Now consider the number M defined as follows: There are two important --- and ultimately contradictory --- facts about the number M: Thus, we have shown that M is divisible by a prime p that is not on the finite list of all prime. And so there must be an infinite number of primes. These two facts imply that M must be divisible by a prime number bigger than pn. Thus, there cannot be a biggest prime. Note that this proof does not provide us with a direct way to generate arbitrarily large primes, although it always generates a number which is divisible by a new prime. Suppose we know only one prime: 2. So, our list of primes is simply p1=2. Then, in the notation of the proof, M=1+2=3. We note that M is prime, so we add 3 to the list. Now, M = 1 +2 *3 = 7. Again, 7 is prime. So we add it to the list. Now, M = 1+2*3*7 = 43: again prime. Continuing in this way one more time, we calculate M = 1+2*3*7*43 = 1807 =13*139. So we see that M is not prime. Viewed another way: note that while 1+2, 1+2*3, 1+2*3*5, 1+2*3*5*7, and 1+2*3*5*7*11 are prime, 1+2*3*5*7*11*13=30031=59*509 is not. Testing for primality. There are a number of simple and sophisticated primality tests. We will consider some simple tests here. In upper-level courses we will consider some faster and more sophisticated methods to test whether a number is prime. Inspection. The most immediate and simple test to eliminate a number n as a prime is to inspect the units digit or the last digit of a number. If the number n ends in an even number 0, 2, 4, 6, 8 we can show that number n cannot be a prime. For example, take n = 87536 = 8753(10) + 6. Since 10 is divisible by 2 and 6 is divisible by 2 then 87536 must be divisible by 2. In general, any even number can be expressed in the form n = a*10 + b, where b = 0, 2, 4, 6, 8. Since 10 is divisible by 2 and b is divisible by 2 then n = a*10 + b is divisible by 2. Consequently, any number n which ends in an even number such as 7777732 or 8896 is divisible by 2 so n is not a prime. In a similar type of argument, if a number n ends in a 5 we can show the number n cannot be a prime. If the last digit of n, call it b, is a 5 we can express n in the form n = a*10 + b, where b = 5. Since 10 is divisible by 5 and b = 5 is divisible by 5 then n = a*10 + b is divisible by 5. Hence, any number n which ends in a 5 such as 93475 is divisible by 5 so n is not a prime. Thus, if a number greater than 5 is a prime it must end with either a 1, 3, 7, or 9. Note that this does not mean all numbers that end in a 1, 3, 7, or 9 are primes. For example, while the numbers 11, 23, 37, 59 are primes, the numbers 21 = 3*7, 33 = 3*11, 27 = 3*9, 39 = 3*13 are not primes. Consequently, if a number ends in a 1, 3, 7, or 9 we have to test further. Trial Division Method. To test if a number n that ends in a 1, 3, 7, or 9 is prime, we could simply try the smallest prime number and try to divide it in n. If that doesn't divide, we would take the next largest prime number and try again etc. Certainly, if we took all primes numbers in this manner that were less than n and we could not divide n then we would be justified in saying n is prime. However, it can be shown that you don't have to take all primes smaller than n to test if n is prime. We can stop earlier by using the Trial Division Method. The justification of the Trial Division Method is if a number n has no divisors less than or equal to formula_41 then n must be a prime. We can show this by contradiction. Let us assume n has no divisors less than or equal to formula_41. If n is not a prime, there must be two numbers a and b such that formula_43. In particular, by our assumption formula_44 and formula_45. But then formula_46. Since a number can not be greater than itself the number n must be a prime. Trial Division Method is a method of primality testing that involves taking a number n and then sequentially dividing it by primes up to formula_41. For example, is 113 prime? formula_48 is approximately 10.63... We only need to test whether 2, 3, 5, 7 divide 113 cleanly (leave no remainder, i.e., the quotient is an integer). So we need not look at any more primes such as 11, 13, 17 etc. less than 113 to test, since 2, 3, 5, 7 does not divide 113 cleanly, 113 is prime. Notice that after rejecting 2 and 3 as a divisor, we next considered the next prime number 5 and not the next number 4. We know not to consider 4 because we know 2 does not divide 113. If 2 cannot divide 113 then certainly 4 cannot because if 4 divided 113 and since 2 divides 4 then 2 would divide 113. So we only use the next cheapest available prime to test not the next consecutive number. If we test 91 we get, So we know since 7 divides 91, 91 is not a prime. Trial division is normally only used for relatively small numbers due to its inefficiency. However this technique has the two advantages that firstly once we have tested a number we know for sure that it is prime and secondly if a number is not prime it also gives us the number's factors. To obtain a few small primes, it may be best to use the Sieve of Eratosthenes than to test each number sequentially using trial division. The Sieve of Eratosthenes method is basically a process of finding primes by elimination. We start by taking a list of consecutive numbers say 1 to 100. Cross out the number 1 because the number is not prime. Take the next least uncrossed off number which is 2 and circle it. Now cross out all multiples of 2 on the list. Next take the next least uncircled number which is 3. Circle the number 3 and cross out all multiples of 3. The next least uncircled number should be 5 since 4 is a multiple of 2 and should have been crossed off. Circle the number 5 and cross out all multiples of 5. The next least uncircled number should be a 7 since 6 is a multiple of 2. Circle the 7 and mark off all multiples of 7. Now the next uncrossed off number should be 11 since 8,9,10 is a multiple of 2, 3, and 2. If we continue in this manner what is left is the circled numbers which are primes. But notice we can actually stop now and circle all the unmarked numbers after crossing off multiples of 7 because of the result that since formula_49 any number less than 100 which is not prime must be divisible by 2, 3, 5, or 7. The Fundamental Theorem of Arithmetic. The Fundamental Theorem of Arithmetic is an important theorem relating to the factorization of numbers. The theorem basically states that every positive integer can be written as the product of prime numbers in a unique way (ignoring reordering the product of prime numbers). In particular, The Fundamental Theorem of Arithmetic means any number such as 1,943,032,663 is either a prime or can be factored into a product of primes. If a number such as 1,943,032,663 can be factored into primes such as 11×13×17×19×23×31×59 it is futile to try to find another different combination of prime numbers that will also give you the same number. To make the theorem work even for the number 1, we think of 1 as being the product of zero prime numbers. More formally, Here are some examples. A proof of the Fundamental Theorem of Arithmetic will be given after Bezout's identity has been established. LCM and GCD. Two characteristics we can determine between two numbers based on their factorizations are the "lowest common multiple", the "LCM" and "greatest common divisor", the "GCD" (also "greatest common factor", "GCF") LCM. The lowest common multiple, or the least common multiple, for two numbers a and b is the smallest number designated by LCM(a,b) that is divisible by both the number a and the number b. We can find LCM(a,b) by finding the prime factorization of a and b and choosing the maximum power for each prime factor. In another words, if the number a factors to formula_50, and the number b factors to formula_51, then LCM(a,b) = formula_52 where formula_53 for "i" = "1 to n". An example, let us see the process on how we find lowest common multiple for 5500 and 450 which happens to be 49500. First, we find the prime factorization for 5500 and 450 which is Notice the different primes we came up for both the number 5500 and the number 450 are 2, 3, 5, and 11. Now let us express 5500 and 450 completely in a product of these primes raised to the appropriate power. The LCM(5500,450) is going to be in the form 2? 3? 5? 11?. All we now have to do is find what the powers of each individual prime will be. So now we compare the power of each prime for 5500 and 450. Let us consider the powers of the first prime 2. In the number 5500, the prime 2 is raised to the second power and in the number 450, prime 2 is raised to the first power. Since the maximum between 2 and 1 for the power of the prime 2 is 2, we use 2 for the power of the prime 2. Now let us consider the powers of the prime 3. In the number 5500, the prime 3 is raised to the zero power and in the number 450 the prime 3 is raised to the second power. Since the maximum between 0 and 2 for the power of the prime 3 is 2, we use 2 for the power of the prime 3. Similarly, let us consider the powers of the next prime 5. In the number 5500, the prime 5 is raised to the third power and in the number 450 the prime 5 is raised to the second power. Since the maximum between 3 and 2 for the power of the prime 5 is 3, we use 3 for the power of the prime 5. Finally, let us consider the powers of the prime 11, the last prime on our list. In the number 5500, the prime 11 is raised to the first power and in the number 450 the prime 11 is raised to the zero power. Since the maximum between 1 and 0 for the power of the prime 11 is 1, we use 1 for the power of the last prime 11. Consequently, the product of our results is LCM(5500,450)=22 32 53 111 = 49500. GCD. The greatest common divisor for two numbers a and b is the biggest number designated by GCD(a,b) that divides both the number a and the number b. In a similar process to finding LCM(a,b), we can find GCD(a,b) by finding the prime factorization of a and b but choosing the minimum power for each prime factor instead. In other words, if the number a factors to formula_50, and the number b factors to formula_51, then GCD(a,b) = formula_52 where formula_57 for "i" = "1 to n". An example, let us see the process on how we find the greatest common divisor for 5500 and 450 which happens to be 50. First, we find the prime factorization for 5500 and 450 which is Notice the different primes we came up for both the number 5500 and the number 450 are 2, 3, 5, and 11. Now let us express 5500 and 450 completely in a product of these primes raised to the appropriate power. The GCD(5500,450) is going to be in the form 2? 3? 5? 11?. All we now have to do is find what the powers of each individual prime will be. So now we compare the power of each prime for 5500 and 450. Let us consider the powers of the first prime 2. In the number 5500, the prime 2 is raised to the second power and in the number 4&lt;/syntaxhighlight&gt; 50, prime 2 is raised to the first power. Since the minimum between 2 and 1 for the power of the prime 2 is 1, we use 1 for the power of the prime 2. Now let us consider the powers of the prime 3. In the number 5500, the prime 3 is raised to the zero power and in the number 450 the prime 3 is raised to the second power. Since the minimum between 0 and 2 for the power of the prime 3 is 0, we use 0 for the power of the prime 3. Similarly, let us consider the powers of the next prime 5. In the number 5500, the prime 5 is raised to the third power and in the number 450 the prime 5 is raised to the second power. Since the minimum between 3 and 2 for the power of the prime 5 is 2, we use 2 for the power of the prime 5. Finally, let us consider the powers of the prime 11, the last prime on our list. In the number 5500, the prime 11 is raised to the first power and in the number 450 the prime 11 is raised to the zero power. Since the minimum between 1 and 0 for the power of the prime 11 is 0, we use 0 for the power of the last prime 11. Consequently, the product of our results is GCD(5500,450)=21 30 52 110 = 50. The Euclidean algorithm. The Euclidean algorithm is such that we can find the gcd of two numbers without finding the factorization*. The Euclidean algorithm consists of only addition and multiplication, and uses the above properties of gcd as its basis. An example. We will see how this works by calculating gcd(458,44) First, divide 458 by 44 and obtain the remainder: Now suppose that a number is a common divisor of 458 and 44. Then it must also be a divisor of 18. To see this, rearrange the above equation to: When this equation is divided by a common divisor of 44 and 458, an integer is obtained on the left, and so must also be obtained on the right. This, by definition, means that the number is also a divisor of 18. By the same reasoning, any common divisor of 18 and 44 is also a divisor of 458. Since all of the common divisors of 458 and 44 are equal to common divisors of 44 and 18, then in particular the greatest common divisors are equal. So we have gcd(458,44)=gcd(44,18) The next step in the algorithm is to divide 44 by 18 and find the remainder. Repeat this process; keep dividing the previous divisor by the previous remainder: Our gcd is the last remainder before the zero, in this case, 2. This is because the reasoning that proved gcd(458,44)=gcd(44,18) applies at every step, so gcd(458,44)=gcd(44,18)=gcd(18,8)=gcd(8,2)=gcd(2,0)=2. The Matrix Method. We can construct a matrix that provides an alternative method for calculating the greatest common divisor. In its general form, the matrix is formula_58 Recall that one way to write the gcd of two numbers is as an integral linear combination. If we are finding the gcd, for example, we could represent it as "as + bt", where "a" and "b" are the two numbers we are comparing, and "s" and "t" are integers. We also know that "b = aq + r" where "r" is the remainder upon division of "b" by "a". After we subtract row 2 from row 1, we get formula_59 If r_2 is nonzero, we must continue the process; this time, subtracting row 1 from row 2. We continue this process until one of the " r's " has been reduced as far as possible. We now have our gcd. The numbers that are in that row, where the 1 and the 0 used to be, represent "t" and "s", respectively. Let us now look at a computational example. formula_60 We see that it would be trivial at this point to go any further; we would just end up with row-2 containing a zero where "a" used to be. So we look at row-1 and remember that the "1" represents our remainder, 1(=t) multiplies "b" and -14(=s) multiplies "a" such that formula_61 This can be checked by the Euclidean algorithm that gcd(7,99)=1. The extended Euclidean algorithm. What happens if we try and reverse the process of the Euclidean algorithm by substituting back? Back-substitution is rather tedious and prone to error, so here is a faster method. Draw up a table with four columns, label these from left to right "q", "r", "u", "v". For convenience label a column "i" representing the step we're currently up to. Place "a" and "b" with the greater of these on top in the column "r", and place 1s and 0s accordingly: formula_62 Now iterate downwards by taking the quotient of "b"/"a" and putting it in the next space in the "q" column, then of "b"-"aq" in the "r" column. To update "u" and "v", take Indeed, you are looking for "u" and "v" such that a"u" + b"v" = gcd (a,b). At some point, gcd (a,b) is in fact the remainder at the ith stage, so you might as well compute ui and vi such that a"u"i + b"v"i = ri, at EACH stage. Deriving the recurrences found above results from these three equations (the second equation is Euclid's algorithm's basic property, the other two are constraints we set to attain our desired goal): The trick is to then appropriately express ri-2. Stop writing when you obtain a 0 in the "r" column. Finding then, gcd(450,44) (this is the same as gcd(44,450) ) formula_63 The bold number is the gcd. Observe (9)×450+(-92)×44=2 Clearly these "u" and "v" are very special. What can we say about the general case? Bezout's identity. In the above case we have 9×450+(-92)×44=gcd(450,44). So the greatest common divisor of 450 and 44 can be written as a linear combination of 450 and 44 with integer coefficients. This turns out to be true of any pair of integers. This result is known as "Bezout's Identity" and can be stated as follows: "Proof" The numbers "u" and "v" can either be obtained using the tabular methods or back-substitution in the Euclidean Algorithm. Proof of the Fundamental Theorem of Arithmetic. One use of Bezout's identity is in a proof of the Fundamental Theorem of Arithmetic. Before this is proven, two other results are needed: Lemma 1: If a prime number, "p", divides a product of two integers, formula_72, then it must divide "a" or "b" (or both). Lemma 2: If a prime number, "p", divides a product of integers, formula_75, then it must divide at least one of the factors. Fundamental Theorem of Arithmetic: Any positive integer, "n", can be expressed as a product of primes. This product is unique up to the order in which the terms appear. Partitioning the Divisors of Products. The Fundamental Theorem of Arithmetic can also be derived from the following lemma: Lemma: Given integers formula_2, formula_3, and formula_93, if formula_93 divides formula_72 (denoted by formula_96), then there exist integers formula_97 and formula_21 such that formula_99 and formula_100 and formula_101. In other words, an integer that divides a product can itself be factored into a product where each factor divides the corresponding factor from formula_72. This means that no new primes are "created" when formula_2 and formula_3 are multiplied together. This Lemma follows from the Fundamental Theorem of Arithmetic and Bezout's identity, but here a more direct proof will be given. Proof: If any of formula_2, formula_3, or formula_93 is formula_108, then the Lemma is trivial. In addition, if any of formula_2, formula_3, or formula_93 is negative, then if the Lemma holds for the absolute values formula_112, formula_113, and formula_114, then it is trivial to show that the Lemma holds for formula_2, formula_3, and formula_93. It will now be assumed that formula_2, formula_3, and formula_93 are all strictly positive integers. Form an formula_121 array formula_122 of integers that has formula_2 columns and formula_3 rows. formula_125 will denote the integer at column formula_19 and row formula_24. Fill the array by sweeping the array row by row starting with row 1, with each row swept starting from column 1. During this "raster" sweep of formula_122, assign values to formula_125 using the following cyclical pattern: formula_130. In essence, formula_125 is the unique integer from the range formula_132 such that formula_133. Since formula_96, it is the case that formula_135 and formula_136. As previously indicated, the "raster sweep" through array formula_122 is a cyclical progression through the entries of formula_122 where column index formula_19 cycles around formula_140, and every time formula_19 transitions from formula_2 to formula_143, row index formula_24 advances by one step around the cycle formula_145. In the image below, the grid formula_122 where formula_147; formula_148; and formula_149 is depicted both explicitly and using a brickwork pattern. Array formula_122 can be endlessly replicated and used to form the infinite array formula_151. For arbitrary integers formula_152 and formula_153, the block of entries in formula_151 formed by columns formula_155 to formula_156 and rows formula_157 to formula_158 is a copy of formula_122. For arbitrary integers formula_19 and formula_24, the entry formula_162 of formula_151 is the unique integer from the range formula_132 such that formula_165. Given any column formula_19 and row formula_24, the entry formula_162 of formula_151 located at formula_170 is the unique integer from formula_171 such that formula_165. Given an arbitrary displacement column displacement formula_173 and row displacement formula_174, the difference formula_175 is separated from formula_176 by a multiple of formula_93. This gives the entries of formula_151 the following symmetries: The columns of formula_151 that contain formula_143 are spaced evenly due to the aforementioned symmetry. Let formula_97 denote the smallest positive integer such that every formula_198 column contains formula_143. The rows of formula_151 eventually repeat (not allowing any column shifts) with a period of formula_21. A row does not appear twice in a single cycle due to the symmetry of formula_151. Row formula_143 is identical to row formula_214 so it must be the case that formula_3 is an integer multiple of the period formula_21: formula_101. It will now be proven that formula_99 by showing that a sub-block of formula_151 that consists of formula_97 columns and formula_21 rows contains every integer from formula_171 exactly once. To clarify notation, given the column indices formula_31 where formula_224, and the row indices formula_32 where formula_226, then formula_227 will denote the sub-block of formula_151 consisting of columns formula_229 through to formula_230, and rows formula_231 through to formula_232. Since a row can be uniquely determined from a single cell and the rows only repeat with a period of formula_21, any block formula_234 will contain exactly formula_21 unique entries. Columns that contain formula_143 occur with a period of formula_97. Due to the symmetry of formula_151, given any integer formula_239, columns that contain formula_240 occur with a period of formula_97. Any block formula_242 will contain exactly formula_243 unique entries. Given any integer formula_239, integer formula_240 will occur in every formula_198 column, and within that column, in every formula_247 cell. Any block formula_242 will contain formula_240. Any block formula_242 will contain every integer from formula_171 exactly once. So therefore formula_99. formula_253. Solving linear modular equations - back to Bezout. Bezout's identity above provides us with the key to solving equations in the form Coprime case - gcd("a", "m") is 1. Consider the case where but with gcd("a", "m")=1 Because of Bezout's identity When we calculate "u", this number is special. Say if we have the equation 4 and 21 are coprime since gcd(4,21)=1. Now 1=4*16+(-3)*(21). Our "u" in this case is 16. Observe now that 4*16=64. 64 (mod 21) = 1. This number "u" is very special - it is known as the "multiplicative inverse". It is the number "u" on multiplication by "a" gives 1 mod "m". Bezout's identity on calculating gcd("a", "m") will always give you the multiplicative inverse of "a" modulo "m". The multiplicative inverse of "a" is often written "a"-1 but note that this does not mean 1/"a" since we have seen in the first sections that we can not always divide in the integers. Note that in Z"p" there is one number "without" a multiplicative inverse - 0. It may be useful to exclude 0 when considering modular arithmetic, so instead of having to say Z"p"\{0} all the time, we merely write Z"p"*. Now since we have the magic multiplicative inverse, our problem becomes relatively easy to solve. 4-1=16 in Z21 and now, on multiplying throughout by 16 (since 4×16=1 because 16 is 4's multiplicative inverse mod 21). 11×16=176 and using a calculator or using the division theorem we obtain which is our solution! Verify - 8×4 = 32 = 11 (mod 21). The general case. Consider the general case where with no restrictions on "a", "b" and "m". Firstly we calculate gcd("a", "m") again to obtain "d". Now "d" is a "divisor" since the "d" in gc"d" means greatest common divisor. So we can now divide "a" and "m" - but what about "b"? Since we have calculated the gcd of "a" and "m" but not "b" we have no guarantees that "d" will divide "b". This then becomes a condition that the equation has no solution. Now we have reduced the problem to the previous coprime case because gcd("a"/"d", "m"/"d")=1 with "d" as above. However we do not have 1 solution any more - this is true because we have reduced the solution to being "x" = "c" (mod "m"/"d") and we must bring the solution back mod "m". This will be come clearer in the examples. Let's work through some examples. Example 1. Solve 4"x" ≡ 3 (mod 20). Firstly, gcd(4, 20) = 4. 4 does not divide 3 and we have no solution. Example 2. Solve 9"x" ≡ 6 (mod 15). gcd(9, 15) = 3 and 3 does divide 6 and we have 3 solutions. Now, divide through by 3 to obtain gcd(3, 5) = 1 = 3 × 2 + -1 × 5 So the inverse of 3 mod 5 is 2. Now we obtain the solution Now in Z15 we must obtain the two extra solutions 9 and 14 mod 15 - 9 mod 5 = 4 and 14 mod 5 = 4. Generally we can say that if we have the solution to the reduced equation "x", the general solution is "x"+("m"/"d")"k" for "k"={0, 1, .., "d"-1}. Chinese Remainder Theorem. Very often congruence relations are required to hold simultaneously. Given positive integer bases formula_254 and formula_255 and arbitrary integers formula_2 and formula_3 where formula_258 and formula_259, a common question is what integers formula_19 satisfy the following congruence relations simultaneously: formula_261 The Chinese Remainder Theorem dictates that when formula_262, for any choice of formula_263 and formula_264 there exists a unique integer formula_265 such that: In essence, when formula_254 and formula_255 are coprime, there is a 1-to-1 correspondence between ordered pairs from formula_271 and the set formula_272. Proof 1. Proof: To begin, observe that formula_273 and formula_274 so it is possible to pair each formula_275 with a unique ordered pair formula_276 and vice-versa. Given any formula_275, integer formula_93 can be reduced modulo formula_254 to get an integer formula_280, and can be reduced modulo formula_3 to get an integer formula_282. Integers formula_2 and formula_3 satisfy: formula_285 It is not obvious that for any choice of formula_286 that there exists a unique formula_275 such that formula_288 Let formula_122 denote an infinite array with two rows indexed by formula_290, and an infinite number of columns indexed by formula_291. formula_125 will denote the entry of formula_122 at column formula_19 and row formula_24. formula_125 is the unique integer from formula_297 where formula_298. Given column indices formula_31 where formula_224, then formula_301 will respectively denote the sub-blocks formed by columns formula_229 through to formula_230 and row sets formula_304. Partition formula_122 into the series of formula_306 blocks formula_307 where block formula_308 is formula_309. Row 1 of block formula_310 will always be the sequence formula_311. Row 2 of block formula_310 can be uniquely determined by its first entry, formula_313 since formula_314. The blocks only differ by row 2, and row 2 for each block is uniquely determined by its first entry formula_313. formula_316 and formula_317 so formula_318. This implies that the first entry of row 2 for the next block can be uniquely determined from the first entry of row 2 for the current block. Since each block is uniquely determined from the previous block, the row 2 pattern for each block will repeat after a regular period of formula_319 blocks: formula_320. Let set formula_321 denote the total range of values attained by formula_322. It is the case that formula_323. Let formula_324 be the minimum positive difference between any two elements from formula_325. The cyclical nature of the elements from formula_325 makes it easy to show that any element formula_327 is congruent to a multiple of formula_324 modulo formula_255. In essence: formula_330. Since formula_324 is the minimum positive difference between any two elements of formula_325, both formula_254 and formula_255 are multiples of formula_324 (in fact, formula_336). Since formula_337, it must be the case that formula_338. This implies that formula_339 and that formula_340. A total of formula_255 blocks are encountered before any repetition occurs in array formula_122, and therefore all formula_343 possible columns occur exactly once in a column period of formula_343 in array formula_122. This establishes the Chinese Remainder Theorem. formula_253 Proof 2. A second (more intuitive) proof can be derived by constructing a mesh to depict the space formula_347. This mesh is a rectangular array of points with formula_254 columns and formula_255 rows. The columns are indexed from left to right by formula_350, and the rows are indexed from bottom to top by formula_351. Most importantly, the mesh will "wrap around" in the horizontal and vertical dimensions. This means that moving to the right from column formula_352 will return you to column formula_108; moving to the left from column formula_108 will send you to column formula_352; moving up from row formula_356 will return you to row formula_108; and moving down from row formula_108 will send you to row formula_356. The mesh has formula_343 points, and the horizontal and vertical coordinates of each point are the remainders from dividing an integer formula_19 by formula_254 and formula_255 respectively. For convenience, given an arbitrary dividend formula_19, formula_365 and formula_366 will denote the remainders after formula_19 is divided by formula_254 and formula_255 respectively. If the dividend formula_19 is incremented by formula_143, then the coordinate formed by the remainders moves to the right by one step, and up by one step, wrapping around if necessary. formula_372 corresponds to the coordinate formula_373. Increasing formula_19 in steps of formula_143 will trace a ray that originates from formula_373 and moves one step to the right and one step up each interation, wrapping around if necessary. The images below give examples of this ray for formula_377 and for formula_378. In the images below, a copy of column 0 and row 0 appears at the right and top of the mesh respectively to illustrate the wrap around property. When formula_378, formula_254 and formula_255 are not coprime and fail to satisfy the conditions of the Chinese remainder theorem. The ray forms diagonal "stripes" in the mesh, and if these stripes are all equally spaced by formula_143, then the ray passes through every point in the mesh exactly once proving that every remainder pair is possible and hence the Chinese remainder theorem. When formula_378, formula_254 and formula_255 are not coprime and fail to satisfy the conditions of the Chinese remainder theorem, hence the ray does not hit every mesh point. The wrap around property of the mesh makes the mesh "symmetric" in both the horizontal and vertical dimensions, which means that if the wrap around "seams" were moved to any column and row where the ray passes through the lower left corner, then the ray is completely unchanged. This requires that the stripes be equally spaced. Let formula_324 denote this equal spacing, and it will be shown that formula_387 and formula_388. The ray passes through row formula_108 every formula_324 steps to the right. The ray passes through formula_373, and the wrap around property implies that moving formula_254 steps to the right returns to this same intersection point. This can only occur if formula_387. By a similar argument formula_388. If formula_254 and formula_255 are coprime, then formula_338, the stripes are evenly spaced by formula_143, every remainder pair is possible, and the Chinese remainder theorem is therefore true. formula_253 

^Lesson 5^ Yes/No Questions. Form the corresponding yes/no question: 1. Ellos tienen hambre. 2. Nosotros estudiamos español. 3. Fernando es alto. Specific Questions. Put in the correct question word: 4. ¿_________ estás? - Estoy bien. 5. ¿__________ es el hombre alto? 6. ¿_________ hermanos tienes? 7. ¿__________ está España? - Está en Europa. 8. ¿_________ es tu cumpleaños? 9. ¿_________estudias español? - Porque es muy interesante. Soluciones a los ejercicios "Solutions to exercices" ^Lesson 5^ 



GCSE Science/Electricity The reason electricity has become so popular over the last 100 years or so is because it is a very good way of moving "energy" from one place to another. When you buy a light bulb is usually rated in watts. The watt is a unit of "power". In this unit you are going to lean some equations rating current voltage, energy and power and apply the equations. Definition of power. Power is defined as the rate of energy flow. Its unit is the watt — one watt is one joule per second. In electrical circuits the power can be found by multiplying the current and voltage together. One way of remembering this is to use an equation triangle like the one below. Memorize the positions of the symbols. "P" is power, "V" is voltage and "I" is current. If you want to know the power, cover its symbol up and what you have left is "I"×"V", or current times voltage. Want to know the voltage? Cover it up and what you have left is "P"÷"I", or power divided by current. Where does the equation "P" equals "VI" come from? From the definition above power is energy "E" per unit time "t". ( how many joules per second) So where does "P" = "IV" come from? To find out we look at what "I" and "V" are defined as. "I" is electrical current. It is the amount of charge "Q" flowing past any point every sec. The voltage also determines the current, so the power is directly proportional to the voltage and the current. We know that from the definition of the volt, if 1 amp flows for 1 second, with a 1 volt potential difference, 1 joule of energy is used up per second, and therefore the power is 1 Watt. Examples. Michael plugs a 100 W light bulb into the mains, what is the resistance of the bulb? To find the resistance we need to know the current, to find the current we use the equation for power.Covering up the "I" in the triangle gives: Now we can use this in the equation for resistance. How much charge flows through this bulb in 15 seconds? We use the definition of current, to work this out. "I" = "C"/"t" So "C" ="It" = (0.435 A)(15 s) = 6.53 coulombs (C) How much electrical energy is converted into heat and light energy in 15 s? For this we use the equation for power in its energy per second form. "P" = "E"/"t" So "E" = "Pt" = (100 W)(15 s) = 1500 joules (J) Q1) Put Ohm's law into a triangle like the one above. Q2) If a current of 3 A flows through a 12 V heater, how much energy will it transfer in half an hour? How much charge will have flowed through the heater in the same time? Q3) In the USA mains voltage is only 120 V. If a 60 W light bulb is connected to the mains in the USA, what current flows through it ? Q4)What current would flow through the same bulb if plugged into the mains in the UK (230 V) and what would be the power? Before we go on to looking at power transmission in the power lines of the national grid, we need to look at how energy is generated. We'll come back to this topic later on, after we look at the next module, which is on the magnetic effects of current. Summary Answers |«Safety | Electromagnetism» 

Cell Biology | ../Parts of the cell/ | ../Organelles/ Endoplasmic Reticulum | Mitochondria and Chloroplasts » The endoplasmic reticulum (endoplasmic="within the cytoplasm", reticulum="little net"; short : ER) is an important organelle in all eukaryotic cells. Prokaryotic organisms do not have organelles and thus do not have an ER. It's base structure and composition is similar to the plasma membrane, though it is an extension of the nuclear membrane. The ER is the site of the translation and folding of and transport of proteins that are to become part of the cell membrane (e.g., transmembrane receptors and other integral membrane proteins) as well as proteins that are to be secreted or "exocytosed" from the cell (e.g., digestive enzymes). The ER consists of an extensive membrane network of tubes and cisternae (sac-like structures). The membrane encloses a space, the cisternal space (or internal lumen) from the cytosol. This space is acting as a gateway. Parts of the ER membrane are continuous with the outer membrane of the nuclear envelope, and the cisternal space of the ER is continuous with the space in between the two layers of the nuclear envelope. Parts of the ER are covered with ribosomes (which assemble amino acids into proteins based on instructions from the nucleus). Their rough appearance under electron microscopy led to their being called rough ER (RER), other parts are free of ribosomes and are called smooth ER (SER). The ribosomes on the surface of the rough ER insert the freshly produced proteins directly into the ER, which processes them and then passes them on to the Golgi apparatus (Fig. 1). Rough and smooth ER differ not only in appearance, but also in function. While the rough ER manufactures and transports proteins destined for membranes and secretion, the smooth ER has functions in several metabolic processes. It takes part in the synthesis of various lipids (e.g., for building membranes) and steroids (e.g., hormones), and also plays an important role in carbohydrate metabolism, detoxification of the cell, and calcium storage. Proteins that are transported by the ER and from there throughout the cell are marked with an address tag that are called a signal sequence. Günter Blobel was awarded the 1999 Nobel Prize in Physiology or Medicine for his discovery of these signal sequences in 1975. The N-terminus (one end) of a polypeptide chain (e.g., a protein) contains a few amino acids that work as an address tag, which are removed when the polypeptide reaches its destination. Proteins that are destined for places outside the ER are packed into transport vesicles and moved along the cytoskeleton towards their destination. Endoplasmic Reticulum | Mitochondria and Chloroplasts » Cell Biology | ../Parts of the cell/ | ../Organelles/ 

The Ablative Case. The ablative case in Latin has 4 main uses: The different uses of the ablative will be dealt progressively. For a summary of all forms of the ablative, please consult the Appendix. Grammar Part 5: The Power of the Ablative Case. Ablative generally indicates position in time and/or space (i.e. when and where). It can also indicate the idea of ways of getting to a location, abstractly or concretely. Ablative of Means. Exercise. How would you translate "I made the toga by hand"? Answer. Answer: "Togam manu feci". In this case, the word "manu" is in the ablative (see fourth declension list) and thus means "by hand." Exercise. I have my wisdom by means of my teacher. Answer. Answer: "Habeo sapientiam magistro." Ablative of Time. How would you say: "I will arrive at the 5th hour." 'at the 5th hour' is indicating position of time. Thus, it can be put into the ablative case, giving: adveniam quinta hora In general, therefore, in order to say "In the morning", "At nine O'clock," or "In the tenth year," use ablative. It is generally used to refer to a specific time in which something has, does, or will occur. Example: I will leave in the night. Hint: Future tense can be looked up in the appendices of this Wikibook! Hint: to leave- discedo, discedere; night- nox, noctis(This is a third declension word!) Answer. Answer: Discedam nocte. Note the simplicity in which Latin translates the six words into simply two. The ending based language completely negates the need for the words "I," "will," "in," and "the." Ablative of Place. "Naves navigabant mari." The ships were sailing on the sea. The ablative is also useful for showing the location of things, in general where you would use the words on, in, or at. There is an exception for the slightly more archaic locative, which is used with the words "domi" (from "domus, domus, f.", home), "ruri" (from "rus, ruris, n.", country [as opposed to city]), and "Romae" (from "Roma, Romae, f.", Rome), as well as with the names of towns, cities and small islands. Latin has its own way of handling prepositions depending on the nouns and their cases in the sentence, including the versatile "in", which can take many different meanings depending upon the case of the object. Ablative with prepositions. Here are a few prepositions that can take the ablative (for a fuller list, see the lesson on adverbs and prepositions in the previous chapter): As a general rule, when motion is implied, use the accusative instead. Example 3. "Servus ex agris venit." Note: "Ager" ("ager, agri, m.", field) must take an ablative suffix to match the preceding preposition, in this case "e"/"ex". Incidentally, both "ager" and "campus" mean "field," but "ager", like its English derivative "agriculture", connotes a farming field, while "campus" (think "camping" or "college campus") means "open field." The "Campus Martius" was a large field in Rome used for military training. The Vocative Case. While you will rarely need to ask Lupus where the bathroom is in Latin, you may find yourself reading either quotes or letters in which a person is being directly addressed. The case it will be in is the vocative. For example, "Hail, Augustus" will appear in Latin as Ave Auguste, and not Ave Augustus. Each declension has its own form of the vocative singular and plural. They are listed in the table below. Furthermore, in all but the second declension, the nominative and vocative are exactly the same! Examples. The basic form of the imperative is created by dropping the "re" off of the infinitive form of the verb, as in: Amare, which becomes Ama; at least in the singular active form, which is all that these exercises require. More can be found about this subject in the chapter on verbs. 

The pronunciation of Church Latin might be easier, for an average European, than Classical Latin. The way introduced here reflect the medieval way. C is pronounced as [ts] or [tS] before e, i, y, ae, or oe G is alway hard, except before e, i, or y, then soft (like j in jut) H is not pronounced. QU might become [k] rather than [kv] or [kw]. 

Chapter 7 Chapter 7. Plant Systematics The Kinds of Plants. The total number of species of plants is tremendous. Any attempt at representing their great diversity requires a system of ordering or arranging the many plant types—hopefully a system that itself contributes to our knowledge and understanding. We could take the straightforward approach of listing all plants alphabetically by their common name, or perhaps by their species name. This index approach would be handy, but would not tell us much about the plants themselves. We would have to know the plants whose names are included for the list itself to have any meaning to us. However, botanists developed a hierarchical system based on the fact—considering a variety of characteristics—that some plants are clearly more similar to each other than to other plants. Arranging all of the world's plants (and animals and minerals) by similarities to each other was an idea first promoted by Carolus Linnaeus. Categorizing is an important process by which we humans gain understanding of the world around us, and something we all do to some degree as part of our observation of things and events that we encounter. In biology, as the concept of evolution was formulated, it became obvious that this concept could be the basis for categorization. If plants that are similar in form are indeed closely related—at least more closely related than plants that are dissimilar in form—then a system of classification could be devised that reflected these relationships. This approach has important implications. Related plants have common properties, a fact that can be exploited in agriculture and other practical botanical fields. Initially, botanists had but one approach: physical examination. The careful examination (and detailed description) of plant structures allowed for arranging each species within a system that placed all more or less similar plants (in certain "important" features) together. This approach is not as easy as it sounds, but played off of and contributed to the expansion of descriptive botany in the 18th and 19th centuries. One problem that became evident is that as species evolved, unrelated plants could come to resemble each other in many respects. After all, form and function are closely related. Within similar habitats (say deserts), species of very distantly related plants might well evolve towards a similar form. Species do not have (or certainly did not have over geological time) unrestricted access to all places on the planet and species distribution is then an important part of interpreting the evolutionary process. As species evolved, they did so within the constraints to dispersion that existed at the time. This fact provides an important clue: unrelated but similar plants are likely to be distributed far from each other on the Earth's surface; and the corollary: plants that have similar structures but have widely separated distributions, may not be so closely related in an evolutionary sense. At this point it is worthwhile to consider some examples. There are many succulent plants, as this form (typically thick, fleshy stems and/or leaves; often reduction or complete loss of leaves) incorporates adaptations necessary for a plant to survive very dry conditions. Non-botanists are tempted to classify all such plants as types of cacti. In fact, cacti evolved in the New World (the Americas), yet there are many succulents (and many plants that resemble cacti) that are not native to the New World, and evolved independently on the African continent. A large group of such plants are known as the euphorbs. 

The Web grew out of a project at CERN, beginning around 1989, where Tim Berners Lee and Robert Cailliau built the prototype system that became the core of what is now the World Wide Web. The original intent of the system was to make it easier to share research papers among colleagues. The original name of the first prototype was Enquire Within Upon Everything, after a famous 19th century reference work of how-tos. Berners-Lee released files describing his idea for the "World Wide Web" onto the Internet on August 6, 1991. 

The World Wide Web (the "Web" or "WWW" for short) is a hypertext system that operates over the Internet. To view the information, you use a software program called a web browser to retrieve pieces of information (called "documents" or "web pages") from web servers (or "web sites") and view them on your screen. You can then follow hyperlinks on the page to other documents or even send information back to the server to interact with it. The act of following hyperlinks is often called "surfing" the web. Looking further at web browsers, a web browser is an application program that accesses the World Wide Web, which then searches for wanted information on the Internet. The first web browser named Mosaic was developed in the early 1990s. The ease of information access provided by web browsers greatly added to the popularity of the Internet. Companies and individual users alike can use a browser to access untold amounts of information, and its as easy to find as clicking a mouse. The four most popular web browsers are Internet Explorer, Chrome, Firefox, and Netscape. Tight competition has caused for continual improvement in the programs and associated technologies. Web browsers are loaded with ease-of-use features and are customizable to an individual user’s preference. URLs, HTTP and HTML. The core functionality of the Web is based on three standards: the "Uniform Resource Locator" (URL), which specifies how each page of information is given a unique "address" at which it can be found; "Hyper Text Transfer Protocol" (HTTP), which specifies how the browser and server send the information to each other; and "Hyper Text Markup Language" (HTML), a method of encoding the information so it can be displayed on a variety of devices. Tim Berners-Lee now heads the World Wide Web Consortium, which develops and maintains these standards and others that enable computers on the Web to effectively store and communicate all kinds of information. Beyond text. The initial "www" program at CERN only displayed text, but later browsers such as Pei Wei's Viola (1992) added the ability to display graphics as well. Marc Andreessen of NCSA released a browser called "Mosaic for X" in 1993 that sparked a tremendous rise in the popularity of the Web among novice users. Andreesen went on to found Mosaic Communications Corporation (now Netscape Communications, a unit of AOL Time Warner). Additional features such as dynamic content, music and animation can be found in modern browsers. Frequently, the technical capability of browsers and servers advances much faster than the standards bodies can keep up with, so it is not uncommon for these newer features to not work properly on all computers, and the web as seen by Netscape is not at all the same as the web seen by Internet Explorer. The ever-improving technical capability of the WWW has enabled the development of real-time web-based services such as webcasts, web radio and live web cams. Java and Javascript. Another significant advance in the technology was Sun Microsystems' Java programming language, which enabled web servers to embed small programs (called applets) directly into the information being served that would run on the user's computer, allowing faster and richer user interaction. The similarly named, but actually quite different, JavaScript is a scripting language developed for Web pages. In conjunction with the Document Object Model (DOM), JavaScript has become a much more powerful language than its creators originally envisaged. Sociological Implications. The exponential growth of the Internet was primarily attributed to the emergence of the web browser Mosaic, followed by another, Netscape Navigator during the mid-1990s. It brought unprecedented attention to the Internet from media, industries, policy makers, and the general public. Eventually, it led to several visions of how our society might change, although some point out that those visions are not unique to the Internet, but repeated with many new technologies (especially information and communications technologies) of various era. Because the web is global in scale, some suggested that it will nurture mutual understanding on a global scale. Publishing web pages. The web is available to individuals outside the mass media. To "publish" a web page, one does not have to go through a publisher or other media institution, and the potential reader is around the globe, some thought. This to some is an opportunity to enhance democracy by giving voices to alternative and minority views. Some others took it as a path to anarchy and unrestrained freedom of expression. Yet others took it as a sign that hierarchically organized society, mass media being a symptomatic part of it, will be replaced by the so-called network society. Also, the hyper-text seemed to promote a non-hierarchical and non-linear way of expression and thinking. Unlike books and documents, hypertext does not have a linear order from the beginning to the end. It is not broken down into the hierarchy of chapters, sections, subsections, etc. This reminded some of the ideas of Marshall McLuhan that new media change people's perception of the world, mentality, and way of thinking. While not unique issue to the web, hypertext in this sense is closely related to the notion of "death of author" and intertextuality in structuralist literary theory. These bold visions are at least not fully realized yet. We can find both supporting and countering aspects of web usage. First, regarding the increased global unity, indeed, many different kinds of information are now available on the web, and for those who wish to know other societies, their cultures, and people, it became easier. When one travels to a foreign country or a remote town, s/he might be able to find some information about the place on the web, especially if the place is in one of the developed countries. Local newspapers, government publications, and other materials are easier to access, and therefore the variety of information obtainable with the same effort may be said to have increased, for the users of the Internet. At the same time, there are some obvious limitations. The web is so far a very text-centered medium, and those who are illiterate cannot make much use of it. Even among the literate, usage of a computer may or may not be easy enough. It has been known during the late 1990s, though with ample exceptions, that web users are dominantly young males in college or with a college degree. Now the trend has been changing and female and elderly are also using the web, level of education and income are related to the web use, some think (See also the Wikipedia article Digital divide). Another significant obstacle is language. Currently, only a limited number of languages are useable on the web, due to software and standard issues, and none would understand all the available languages. These factors would challenge the notion that the World Wide Web will bring unity to the world. Second, the increased opportunity to individuals is certainly observable in the countless personal pages, as well as other groups such as families, small shops, which are not among those who publish materials. The emergence of free web hosting services is perhaps an important factor in bringing this possibility into reality. The activities of alternative media expanded into the web as well. Yet not a small part of those pages seem to be either prematurely abandoned or one-time practice. Very few of those pages, even when they are well-developed, are popular. When it comes to the expression of ideas and provision of information, it seems that the major media organizations and those companies who became major organizations through their online operations are still favored by the dominant majority. Besides, the Web is not necessarily a tool for political self-education and deliberation. The most popular uses of the Web include searching and downloading pornography, which perhaps has a very limited effect in improving democracy. The most intensively accessed web pages include the document detailing the former U.S. president Bill Clinton's sexual misconduct with Monica Lewinsky, as well as the lingerie fashion show by Victoria's Secret. In sum, both in terms of writers and readers, the Web is not popularly used for democracy. While this is not enough to categorically reject the possibility of the Web as a tool for democracy, the effect so far seems to be smaller than some of the expectations for a quite simple reason, lack of interest and popularity. Anarchistic freedom of expression may be enjoyed by some, but many web hosting companies have developed their acceptable use policy over time, sometimes prohibiting some sensitive and potentially illegal expressions. And again, those expressions may not reach great many. The web is still largely a hierarchical place, some may argue. Third, regarding the non-linear and non-hierarchical structure of the Web, the effect of those on people's perception and psychology are still largely unknown. Some argue that our culture is changing to that of postmodernity, which is closely related to a non-linear and non-hierarchical way of thinking, being, and even social organization. Yet the counter-evidence is available as well. Among the most notable would be the existence of web directories and search engines. Those sites often provide navigations to the most popular sites to visitors. Besides, it is quite obvious that many web sites are organized according to a simple hierarchy, having the "home page" at the top. At least the present state of the Web and web users seem to suggest the change has not been as great as envisioned by some. 

^Lesson 5^ Rellena los espacios en blanco "Fill in the blank". 1. Manuel tiene una bicicleta. Su bicicleta es azul. "Manuel no has a bike. His bike is blue". 2. Mis clases son a las nueve de la mañana. "My classes are at 9:00 am". 3. Nosotros visitamos a nuestros abuelos. "We visit our grandparents". 4. Los mexicanos tienen su día de independencia el 16 de septiembre. "The Mexicans have their independence day on September 16" 5. ¿Pueden mostrarme tus casa? "Can you (plural) show me your house?" 6. La clase de Señor Ford es fantástica. Su clase es fantástica. "Mr. Ford's (his) class is fantastic." Soluciones a los ejercicios ^Lesson 5^ 

Electric currents produce magnetic fields and moving magnetic fields produce electric currents. This module looks that this property of electric and magnetic fields in detail. You will learn about the magnetic fields generated by a wire and a solenoid, then go on to look at motors generators and transformers. Review of Magnetism. From work in lower school you should already be familiar with the basic properties of magnets. They are: Notice the similarity between magnetism and electricity. Electricity attracts neutral materials (recall a rubbed balloon picking up small bits of paper). Electricity comes in two forms:positive and negative. Like charges repel, unlike charges attract. Physicists long ago realised that these similarities were important. It is now realised that electricity and magnetism are closely related forces. They are really two different aspects of the same force! We call that force electromagnetism. This next section of the electricity module is probably the hardest. But it is also the most interesting. Work the problems in the text as you go along and you "will" be able to cope. 

Appendix A: Introduction to reactions &lt;br&gt; Appendix B: Index of reactions &lt;br&gt; Appendix C: Introduction to functional groups Help organize the book structure. Compare this book to these college OChem textbooks: If you think you can help, check out the to do list of the authors over here - /To-Do List/ 

This General Science book is aimed at GCSE students rather than university students. Contents by Topic. /Introduction/ Contents by Modules. The GCSE exam can be modular or combined. It is sensible however to divide the subject into modules despite the exam system used. Although this page is set out using the English system, students from other nations will still find much of the material relevant. Coursework. Common coursework experiments.  __NOEDITSECTION__ 

Question 5. A student rubs a polythene strip with fur and suspends it from a clamp stand. She then rubs another polythene strip with fur and brings it up to the first strip. The two strips repel each other. Look at the following four statements: Which two of the above statements are false? Question 6. A student decided to silver plate his/her mother's forks. S/he set up the apparatus above. S/he set the power pack and variable resistor so that the ammeter read 0.1 A. S/he allowed the fork to remain in the electrolysis apparatus for 10 min. Once s/he removed the fork s/he tested it by trying to scratch the silver off. S/he found that the layer of silver was too thin and decided that s/he would have to take steps to make the layer a lot thicker. Which one of the following steps would not produce a thicker layer of silver ? Question 7. What is the maximum power that an appliance should have if it is connect to a 230V supply by a 5A cable ? Question 8. The graph above shows how the potential difference across a conductor varied with the current flowing through it. Which of the following statements is true ? Question 9. Two students, Jane and John are working together to make an electromagnet. John coils some plastic coated wire around a pencil and attaches the two ends to a 1.5V cell. He finds the electromagnet that he has made will deflect a compass placed near it but will not pick up a paperclip. Jane suggest that using an iron nail rather than a pencil will make the magnet stronger. Is she correct, and why? Question 10. Jane and John have made an electromagnet using a soft iron core with 10 windings of insulated copper wire and a voltage of 1.5V. They found that they could pick up 3 paper clips with this magnet if they were careful. Which of the following changes would "not" result in more paper clips being picked up? Question 11. A photocopying machine creates a picture by making use of static electricity. Put the following statements about how the process works in the correct order. (Don't see your order in among the options? Check again because the correct order is certainly in there). Question 12. A wire is connected to a power pack and placed in the magnetic field of two bar magnets. When the power pack is switched on the wire jumps upwards. What effect will increasing the current in the wire have ? Question 13. This is the same setup as in question 12. What can be done to make the wire move downwards ? 

A kettle is fitted with a 13A fuse which one of the following statements is "correct". 

^Lesson 5^ Llena los blancos. "Fill in the blank". 1. Mario es __________ alto como Luigi. 2. El Sr. Perez tiene __________ perros como Sr. Gonzales. 3. Serena es __________ fuerte que Venus. 4. Mi carro cuesta __________ dinero como tu carro. 5. Los niños juegan __________ horas que los mayores. 6. México es __________ hermoso como España. 7. El abuelo tiene __________ pelo como su hijo. 8. Los gatos son __________ atrevidos como los tigres. 9. Argentina tiene __________ personas como Colombia. Soluciones a los ejercicios "Solutions to exercices" ^Lesson 5^ 

Sorry this is the wrong answer The current is determined by the voltage and the resistance of the kettle element. The fuse can't contol the current. 

Sorry this is the wrong answer. Although the circuit symbol for a fuse looks a bit like a resistor, inside a fuse is just a piece of wire. It can't keep the current down to below 13A in this way. 

Well done! This is the correct answer! Next question 

When a current flows down a wire it creates a magnetic field. To see this, place a small plotting compass near the wire and turn the current on. The compass needle will deflect. The shape of the magnetic field around a wire is circular. Look at the diagram on the right. The wire is coming straight out of the screen so you only see it's cross section (the red circle) The plotting compasses show how the field wraps around the wire. Creating an electromagnet. Solenoids are made by coiling wire. The magnetic field of a solenoid looks like the field of an ordinary bar magnet. A solenoid makes a pretty weak bar magnet on its own, but if a piece of iron is put inside the solenoid, the field becomes much much stronger. Try the following experiment: You will find that the solenoid with the iron core is much much stronger than the one with the air core. Q1) Why is it important to use plastic-coated wire? Making the field stronger. To make the field stronger we can: Reversing the field. To make the north pole and south pole swap positions we can: Speaking of north and south poles, there is a little trick to help find out which end is the north pole and which is the south. Look down the solenoid and work out which way the current is flowing. Remember, current flows from the positive to the negative terminal on the power pack. If the current is flowing clockwise, the end will be a south pole. If it is going counterclockwise, the end will be a north pole. The easy way to remember this is to put arrows on the end of a capital N and S like this: If you do not understand how to determine the pole using the clockwise and anti-clockwise way you can use the right hand grip rule 

 __NOEDITSECTION__ The Constitution and Government of the United States of America. Part IV- American Government: Theory and Analysis. Part V- Appendices 

Introduction. The United States Constitution divides government into three separate and distinct branches: the Executive, Legislative and Judicial branches. The concept of separate branches with distinct powers is known as "separation of powers." That doctrine arose from the writings of several European philosophers. The Englishman John Locke first pioneered the idea, but he only suggested a separation between the executive and legislative. The Frenchman Charles-Louis de Secondat, Baron de Montesquieu, added the judicial branch. Each branch is theoretically equal to each of the others. The branches check each others powers and use a system known as checks and balances. Thus, no branch can gain too much power and influence, thus reducing the opportunity for tyrannical government. The Preamble to the American Constitution sets out these aims in the general statement: "We the people of the United States, in order to form a more perfect union, establish justice, insure domestic tranquility, provide for the common defense, promote the general welfare, and secure the blessings of liberty to ourselves and our posterity, do ordain and establish this constitution for the United States of America". The Legislative. The Congress is the Legislative Branch. Its main function is to make laws. It also oversees the execution of these laws, and checks various executive and judicial powers. The Congress is "bicameral" - it is composed of two houses. One house is the House of Representatives and the other is the Senate. The House of Representatives is currently composed of four hundred and thirty-five members. Each of the fifty states is allocated one or more representatives based on its population which is calculated on a decennial basis (once in ten years) . Each state is guaranteed at least one representative. A state that is allocated more than one representative divides itself, as state procedures dictate, into a number of districts equal to the number of representatives to which it is entitled. The people of each district vote to elect one representative to Congress (States that have only one representative allocated choose "at-large" representatives - the state votes as one entire district). The District of Columbia and a number of U.S. territories have been permitted to elect delegates to the House of Representatives. These delegates may participate in debates, and sit and vote in committee, but are not allowed to vote in the full House. Every House member faces re-election in an even-numbered year and is elected to a two-year term. The House is presided over by a Speaker, who is directly elected by the members of the House. The Senate is the upper house of the legislative branch of the United States and possesses one hundred members which is considerably less than the four hundred and thirty-five members of the House of Representatives. Each state chooses two senators, regardless of that state's population. The Constitution originally dictated that a state's senators were to be chosen by the state's legislature; after the Seventeenth Amendment was ratified in 1913, senators were elected directly by the state's population. In contrast to the House's two-year terms, Senators are elected to a six-year stint in office. In addition, only one-third of the Senate stands for election during an even year. These differences between the two houses were deliberately put into place by the Founding Fathers; the Senate was intended to be a more stable, austere body, whereas the House would be more responsive to the people's will. The Vice-President is President of the Senate, but he/she only votes if there is a tie. The Senate also chooses a President Pro Tempore to preside in the Vice-President's absence (though, in practice, most of the time, senators from the majority take turns presiding for short periods). The Senate and the House are both required to approve legislation before it becomes a law. The two houses are equal in legislative power, but revenue bills (bills relating to taxation) may only originate in the House. However, as with any other bill, the Senate's approval is still required, and the Senate may amend such bills. The Senate holds additional powers relating to treaties and the appointments of executive and judicial officials. This power is known as "advice and consent." The Senate's advice and consent is required for the President to appoint judges and many executive officers, and also to ratify treaties. To grant advice and consent on treaties, two-thirds of the Senators must concur (agree). While most votes require a simple majority to pass, it sometimes takes three-fifths of senators to bring a bill to a vote. This is because Senate rules hold that a bill cannot be voted on as long as it is being debated--and there is no limit on how long a senator may debate a bill. Senators sometimes use this rule to filibuster a bill--that is, continue debating a bill endlessly so that it cannot be voted on. The only way to end a filibuster is for three-fifths of all Senators to vote for a cloture resolution, which ends all debate and brings the bill up for voting. Use of the filibuster tends to be controversial. Whichever party is in the majority tends to call its use "obstructionism," while the other side sees it as an important check on the majority. The House has the sole power to impeach federal executive and judicial officers. According to the Constitution, officers may be impeached for "treason, bribery, or other high crimes and misdemeanors.” The Senate has the sole power to try all such impeachments, a two-thirds vote being required for conviction. The Constitution requires that any individual convicted by the Senate to be removed from office. The Senate also has the power to bar that individual from further federal office. The Senate may not impose any further punishment, although the parties are still subject to trial in the courts. As the Vice-President (being next-in-line to the Presidency) would have an obvious conflict of interest in presiding at a trial of the President, in such cases, the Chief Justice presides. Interestingly, no similar provision prevents the Vice-President from presiding at his or her own trial. The Executive. The President, Vice President, and other executive officials make up the Executive Branch. The main function of this branch is to execute the laws created by Congress. The President and the Vice-President are chosen by the Electoral College, a body of people elected for the purpose of electing the President. One may wait to consider the Electoral College in further detail. The President appoints several "Secretaries" to head executive departments. An executive department is a body covering a broad topic of law- examples include the Department of Agriculture and the Department of Justice. The several secretaries (in the case of the Justice Department, the Attorney General) serve as advisors to the President and also as the chief officers of their own departments. This group of advisors is collectively known as the President's "cabinet". The President nominates these Secretaries, as well as other important federal officials, and the Senate "advise and consents" to them . The Judicial. The Supreme Court and the lower courts compose the Judicial Branch. The judiciary must interpret the laws of the United States. In the course of such interpretations, the courts may find that a law violates the constitution. If so, the court declares the law unconstitutional. Thus, the judiciary also has a role in determining the law of the land. The judges of federal courts are nominated by the President and advised and consented to by the Senate. The number of judges and the exact structure of the courts is set by law, and not by the Constitution. The Legislative Process: How A Bill Becomes A Law. After both houses of Congress pass a bill, perhaps observing the different rules and procedures in each house, but with the exact same final text, the bill is submitted to the President. Immediately, a ten-day clock for the president to act in starts to tick. Sundays are excluded in this calculation. Once he receives the bill, the President has many options. The outcome of the process depends on the route taken by him. Checks and Balances. In order to prevent any branch of government from becoming too powerful, the Framers of the Constitution created a system of checks and balances. Each branch of government has checks on the others, while it is itself also checked. The complex system can be outlined as follows: Checks of the Legislative "Checks on the Executive" "Checks on the Judicial" "Internal Checks" Checks of the Executive "Checks on the Legislative" "Checks on the Judicial" Checks of the Judicial "Checks on the Legislative" "Checks on the Executive" 

Welcome to the Wikibook of ALGEBRA Preface. Algebra is the branch of mathematics concerning the study of the rules of operations and relations, and the constructions and concepts arising from them, including terms, polynomials, equations and algebraic structures. You've been asked to buy peanuts for you and your friends at a football game. You've collected $12.50. One bag is $2.75. You want to know how many bags you can buy. This is an algebra problem! Related problems are how much money will be left over, and what should you buy or in what proportions should you return the extra money to your friends. Algebra helps us to predict things that we don't yet know and to determine the relationships between the things we do. Algebra is a powerful and rich branch of mathematics that is useful in everyday life as well as business, engineering, and other technical fields. Contents. Unit 1: Numbers, Variables and Relationships. How can we use numbers and variables to find out unknown information? Chapter 1: Elementary Arithmetic&lt;br&gt; Chapter 2: An Introduction to Algebra&lt;br&gt; Chapter 3: Solving Equations&lt;br&gt; Chapter 4: Inequalities&lt;br&gt; Unit 2: An Introduction To Graphing. Math is a method of solving problems. You take information you know, and by manipulating it using mathematical principles, you can find information you don't know. Functions are the mathematical framework for solving problems. They have parameters, rules, and ways of being solved. This section will introduce you to functions and how to use them. Unit 7: Advanced Algebra. Now that you have completed your mathematical journey, you may be wondering where to head over next. This unit provides six optional chapters that can be done in any order you please. The six chapters are meant to be independent from one another, and any of the chapters can be skipped entirely to put more focus on the topics that are of greater interest to you. 



Quadratic equations. Up to now you have only dealt with equations and expressions involving just x; in this section we'll move onto solving things which have formula_1 in them. All quadratic equations can be arranged in the form formula_2, and "a","b","c" are all constants. Now let's look at some examples: Examples: Rearrange the following equations in the form formula_3: Solution for (1): Solution for (2): Factorization. Factorization is the most common way to solve quadratic equations. Let us consider again the first example above: formula_11 We have already simplified the equation into Now, we want to factorize the equation - that is to say, get it into a form such as: Look at the number term c. In this example, it is -3. Now, if we are lucky, the numbers "something" and "something else" will turn out to be nice whole numbers, so let's think of two numbers that will multiply together to give -3. Either 3 and -1, or -3 and 1. But we also need to get the x term correct (here, b=2). In fact, we need our two factors of c to add together to make b. And (3)+(-1)=2. So, we have found our 'somethings': they are 3 and -1. Let's fill them in. Just to check, we can multiply out the brackets to check we have what we started with: Now, we know that in an equation the left side is always equal to the right side. And in this case the right side of the equation is 0, so from that we can conclude the term formula_16 must equal to zero as well. And that means that either formula_17 or formula_18 must equal zero. (Not convinced? Remember (x+3) and (x-1) are just numbers. Can you find two non-zero numbers which multiply to make zero?) Let's write that algebraically: Thus, there are two different solutions to the same equation! This is the case for all quadratic equations. We say that this quadratic equation has "two distinct and real roots". With practice, you will often be able to write down the equation in factorised form almost immediately. Here is another example, in this case the x easily factorises out: Completing the square. Sometimes the roots (solutions) of a quadratic equation cannot be easily obtained by factorisation. In such cases, we have to solve the equation by completing the square, or using the quadratic formula (see below). In order to complete the square, we need to rewrite the given equation in the form formula_21. Now here is an example: In general, we get Note that when we reach the stage of taking the square root of both sides of the equation, we might have a negative left-hand side. In this case, the roots will be complex. If you have not yet learned about complex numbers, it is possible to simply state that the equation "has no real roots". Quadratic Formula. The quadratic formula is a special generalization of completing the square that allows the two roots of a quadratic equation to be obtained by simple substitution. It can be used to solve any quadratic equation and is very quick to work out on a calculator. Complete the square: Simplify: Which equals 4y to the 19th power. which is the desired form of the quadratic formula. Hence, given that a quadratic is in the form formula_32, the two roots are: The quantity formula_34 in the equation, known as the discriminant, is an indication of the solubility and nature of the roots: Weda's Theorem. If the quadratic equation formula_35 has two real roots formula_36 and formula_37, then formula_38 This is because formula_39 and formula_40. By simply adding or multiplying the two roots we will get the above two equations. This is called Weda's Theorem. Using Weda's Theorem we can find the second root of a given quadratic equation without solving the equation. Example: Given that one of the real roots of the equation formula_41 is 2, find the other root without solving the equation. Solution: formula_42 We can also determine the signs of two roots by applying the following rules: Another problem involving Weda's Theorem: Example: For the equation formula_47, given that the sum of squares of roots is formula_48, find the value of formula_49. Solution: formula_50 Pythagorean Theorem ____________________ a^+b^=c^ Solving simultaneous linear and nonlinear equations. In previous chapters you have already learned how to solve simultaneous linear equations. Now we will learn how to solve a system of simultaneous linear and non-linear equations with two unknowns. It is usually done by substitution method. Example: Solve the following simultaneous equations: formula_51 Solution: formula_52 ∴ x=-1 and y=1, or x=-2 and y=0.  

Introduction. The United States is exactly that--a Union of states. Each state has its own individual powers. However, that does not mean that the states have power to legislate on all matters. The Constitution of the United States spells out the powers of the federal government and of the "several states." The Union government (known as the federal government) has its own fields of legislation, and if federal legislation conflicts with the state laws, the federal legislation prevails. If this occurs, the state must defer to the federal government. The alternative, that any state may at any time leave the Union and thus be free from Union interference in the state's internal affairs, was tried during the American Civil War. Types of Federalism. There are two types of federal systems. The first, dual federalism, holds that the Union and the state are equal; under this view of federalism, the Union government only has the powers expressly granted to it, while the states retain all other powers. The second view, cooperative federalism, states that the federal government is definitively superior to the state government, and the federal government should stretch its powers as far as possible. The US is a Union that does not completely fit either definition. The type of federalism in effect really depends on who is in power at the time. In any case, the Constitution prevents stretching federalism too far in either extreme. All powers retained by the states are known as reserved powers. Those specifically granted only to the Union government are known as enumerated powers. Finally, matters over which both the Union and the state governments have control are known as concurrent powers. The Tenth Amendment provides that the Union government has only the powers expressly designated to it by the Constitution, and the states control all other matters. Enumerated powers relate to the following: Reserved Powers. As has been mentioned, the states have control over all matters not controlled by the Union government. Some of these matters include: Concurrent Powers. Several powers belong concurrently to the Union and the state governments. These include: 

Introduction. The Constitution provides for many things other than the Union and member state government operation. Supreme Law. The Constitution calls itself, the United States laws, and treaties, the Supreme Law of the land. No state law may conflict with the Supreme Law. If they do, the conflicting parts are ruled void. However the supremacy of federal law over state law only applies if the federal government is acting in pursuance to its constitutionally granted powers. Also, the Constitution is itself superior to federal laws and treaties. If either of these, or any Executive order made by the President, or a state law, or other regulation similar to law, conflicts with the constitution, the Courts may declare it unconstitutional, and the law or regulation becomes void. Note that this is not explicitly provided for by the Constitution. Chief Justice John Marshall originally used judicial review- the power to declare laws unconstitutional- in the early 1800's when he decided the case "Marbury v. Madison". Oaths. The Union and member state officers take an oath to support and defend the Constitution. Only the Union President's oath is explicitly spelled out- the form of the other oaths is left to the appropriate government. The Constitution, while mentioning oaths, specifically prohibits the United States or a member state from requiring a religious oath or observance of a certain religion in order to qualify for any office. Amendment. The Constitution is not an unchangeable document. Changes to it are known as amendments. By tradition, an amendment does not strike out words and insert others. In order to leave the original Constitution whole, Amendments are added as separate Articles at the end of the Constitution. There are two ways in which an amendment can be proposed, and two ways in which it can be ratified, or approved. First, the Amendment can be proposed by Congress. For this to occur, two-thirds of the House of Representatives and two-thirds of the Senate must vote for the Amendment. Second, an Amendment can be proposed by a Constitutional Convention. If two-thirds of the member states make an application to Congress, Congress must call this Convention, which then proposes Amendments. However, a Constitutional Convention has never been called. Regardless of the way in which the Amendment is proposed, it must be ratified by three-fourths of the member states. The first manner in which ratification by a member state may occur is through the legislature. Secondly, a state convention can be called to ratify an amendment. The second method has only been used once. Note that each member state cannot decide which method it wishes to use. Whoever proposed the Amendment- the Congress or a Union Convention- will decide if legislatures or state conventions will ratify. There have been twenty-seven Constitutional Amendments to date. The first ten guarantee certain already existing rights to the people and member states- they are therefore known as the Bill of Rights. 

First Amendment. "Congress shall make no law respecting an establishment of religion, or prohibiting the free exercise thereof; or abridging the freedom of speech, or of the press; or the right of the people peaceably to assemble, and to petition the government for a redress of grievances." The First Amendment is the best-known amendment to most Americans. It provides for and guarantees certain fundamental and basic human rights, such as: 1. Freedom of religion. The amendment states: "Congress shall make no law respecting an establishment of religion, or prohibiting the free exercise thereof." This means that Congress cannot make one religion official, or require people to worship in a certain way. It also forbids Congress, Federal and State courts, and State legislatures from prohibiting any one religion from being practiced. 2. Freedom of speech and of the press. The amendment also guarantees that Americans have freedom of speech, and forbids both the Federal Government and State Governments from punishing a person for expressing their views. It also calls for the freedom of the press and bars the Government from interfering with/censoring the press. Therefor the government cannot control the press (i.e. controlling what the press publishes or punishing the same for publishing a particular article or airing on television or radio a particular segment/video.) 3. Freedom of assembly and petition. The amendment guarantees that the people have the right to assemble in a peaceful manner and to consult for their common good. This has also be interpreted by the courts to also guarantee freedom of association. It also states that Americans may petition the Federal Government for a redress of grievances. Second Amendment. "A well regulated militia, being necessary to the security of a free state, the right of the people to keep and bear arms, shall not be infringed." The second amendment guarantees that individual Americans have the right to keep and bear firearms. The amendment is subject to many controversies about whether it applies to only people in the armed forces or everyday citizens. The Supreme Court has ruled however in "District of Columbia vs Heller (2008)" that the amendment applies to everyday citizens. A majority of gun laws and regulation in the USA are at State level. Third Amendment. "No soldier shall, in time of peace be quartered in any house, without the consent of the owner, nor in time of war, but in a manner to be prescribed by law." Soldiers, according to the Third, may not during peacetime be quartered or housed in a residence without the consent of the owner. Furthermore, even in wartime, a law is required to force people to quarter soldiers. The Supreme Court and others have also attributed a right to personal privacy within this amendment. Fourth Amendment. The fourth amendment of the United States Constitution reads as follows: The right of the people to be secure in their persons, houses, papers, and effects, against unreasonable searches and seizures, shall not be violated, and no Warrants shall issue, but upon probable cause, supported by Oath or affirmation, and particularly describing the place to be searched, and the persons or things to be seized. The fourth amendment can be broken into two distinct parts. The first part provides protection against unreasonable searches and seizures, although unreasonable has been defined in a myriad of different ways. However, the framers did not qualify how a violation of this amendment was to be punished. Case law has afforded the exclusionary rule to ensure that evidence improperly collected is excluded from trial. (See "Wolf v. Colorado", "Mapp v. Ohio", and "United States v. Leon" for more information on the progression of the exclusionary rule.) This gives law enforcement officers the incentive to respect the amendment. The second part of the amendment provides for the proper issue of warrants. However, it is important to note that it does not require that warrants be obtained prior to a search or seizure, only that searches and seizures are reasonable. When warrants are issued, their must be probable cause. Probable cause is tested using the "totality of circumstances" test as defined in "Illinois v. Gates". (See "Spinelli v. United States" and "Illinois v. Gates" for the progression on probable cause tests. Also see "Maryland v. Garison" and "Richards v. Wisconsin" for more on search warrants.) Interesting Search and Seizure Cases Fifth Amendment. "No person shall be held to answer for a capital, or otherwise infamous crime, unless on a presentment or indictment of a grand jury, except in cases arising in the land or naval forces, or in the militia, when in actual service in time of war or public danger; nor shall any person be subject for the same offense to be twice put in jeopardy of life or limb; nor shall be compelled in any criminal case to be a witness against himself, nor be deprived of life, liberty, or property, without due process of law; nor shall private property be taken for public use, without just compensation." While the Fourth protects one against the police, the Fifth guarantees one's rights relating to actual charges of criminal activity. Firstly, the Fifth requires that a Grand Jury (a special body of lay-people) indict, or accuse, a person of a serious crime before he may be tried for it. The sole exception relates to the military, where Grand Juries are not used. 2: the Amendment prohibits "double jeopardy." Once a person has been tried and acquitted (found not guilty), the prosecutor cannot retry him, even if new evidence is discovered. However, the prohibition against double jeopardy does not extend to mistrials. If a jury cannot reach a verdict and the judge declares a mistrial, then the case may be retried. 3: the Amendment provides that a person cannot be compelled to testify against themselves. A person may "plead the Fifth" to avoid self-incrimination. However, the government may grant immunity, or formal protection from prosecution, to a person. Then, the person cannot be punished for what crimes he might have committed, and can be forced to give evidence. 4: the Fifth requires due process to be used. While due process is not defined in the Constitution, one may conclude that it involves fair trials with impartial juries and judges. 5: the Fifth prohibits government from taking away a person's property unless it provides fair compensation. Sixth Amendment. In all criminal prosecutions, the accused shall enjoy the right to a speedy and public trial, by an impartial jury of the State and district wherein the crime shall have been committed, which district shall have been previously ascertained by law, and to be informed of the nature and cause of the accusation; to be confronted with the witnesses against him; to have compulsory process for obtaining witnesses in his favor, and to have the Assistance of Counsel for his defence. The sixth amendment, like the fifth amendment, also guarantees and lays out a person's rights when that person had been accused of committing a crime. Firstly, a person is entitled to a "speedy and public trial by an impartial jury." This means a person cannot be arrested/detained for an unreasonable amount of time before that person goes to trial, and the trial cannot be private, unless the person accused waivers this right and ask that the trial be private. The Supreme Court has also ruled that a trial can be held in private, if by having excess publicity would serve to undermine the accused person's right to due process. Also an accused person has the right to trial by jury, and the trial must held in the State or district in which the crime was alleged to be committed. Secondly, an accused person has the right to be informed of the "nature and cause of the accusation" against him. Therefore, an indictment must allege all the ingredients of the crime to such a degree of precision that it would allow the accused person to assert double jeopardy if the same charges are brought up in subsequent prosecution. Thirdly, an accused person has the right to confront/cross-examine the wittiness against him, this also includes the right to cross-examine physical evidence that the prosecution will use against the accused person during the trial. An accused person also has the right to call witness to testify in his favor. If the witness refuses to testify, the court may, at the accused's request, order the witness to do so. However, in some cases the court may refuse to permit a defense witness to testify. For example, if a defense lawyer fails to notify the prosecution of the identity of a witness to gain a tactical advantage, that witness may be precluded from testifying Fifthly, an accused person has the right to be represented by counsel (i.e. a lawyer) and if the accused person does not have a lawyer, and cannot afford one, the court must, upon the accused's request, appoint a lawyer to represent him. If the accused person is non compos mentis and incapable adequately of making his own defense, the court must, with or without the accused's request, appoint a lawyer to represent him. Seventh Amendment. "In suits at common law, where the value in controversy shall exceed twenty dollars, the right of trial by jury shall be preserved, and no fact tried by a jury, shall be otherwise reexamined in any court of the United States, than according to the rules of the common law. The seventh amendment guarantees the right to trial by jury in common law (i.e. civil) suits. Eighth Amendment. "Excessive bail shall not be required, nor excessive fines imposed, nor cruel and unusual punishments inflicted." The eighth amendment prohibits excessive bail from being required. This means that bail must appropriately match the crime in which the arrested/detained person is accused of committing. Therefor a person arrested for littering cannot be put on bail for 50,000 dollars, as that would be excessive for that respective offense. However bail may be denied if a person is arrested on murder charges. It also prohibits excessive fines from being imposed as punishment for crime and again must appropriately match the respective crime. Lastly, it forbids cruel and unusual punishments. Ninth Amendment. The enumeration in the Constitution, of certain rights, shall not be construed to deny or disparage others retained by the people. The ninth amendment states that Americans have more rights than what are listed in the Bill of Rights and that all because a right is not listed in the same, does not mean Americans do not posses that right. Tenth Amendment. "The powers not delegated to the United States by the Constitution, nor prohibited by it to the states, are reserved to the states respectively, or to the people. The Tenth Amendment prevents the Union government from regulating a matter it is not constitutionally entitled to regulate. The Union government, therefore, cannot infringe upon a state's reserved powers. 

The Preamble. We the People of the United States, in Order to form a more perfect Union, establish Justice, insure domestic Tranquility, provide for the common defence, promote the general Welfare, and secure the Blessings of Liberty to ourselves and our Posterity, do ordain and establish this Constitution for the United States of America. Article I. Section 1. All legislative Powers herein granted shall be vested in a Congress of the United States, which shall consist of a Senate and House of Representatives. Section 2. The House of Representatives shall be composed of Members chosen every second Year by the People of the several States, and the Electors in each State shall have the Qualifications requisite for Electors of the most numerous Branch of the State Legislature. No Person shall be a Representative who shall not have attained to the Age of twenty five Years, and been seven Years a Citizen of the United States, and who shall not, when elected, be an Inhabitant of that State in which he shall be chosen. Representatives and direct Taxes shall be apportioned among the several States which may be included within this Union, according to their respective Numbers, which shall be determined by adding to the whole Number of free Persons, including those bound to Service for a Term of Years, and excluding Indians not taxed, three fifths of all other Persons. [Superseded by Amendment XIV, section 2] The actual Enumeration shall be made within three Years after the first Meeting of the Congress of the United States, and within every subsequent Term of ten Years, in such Manner as they shall by Law direct. The Number of Representatives shall not exceed one for every thirty Thousand, but each State shall have at Least one Representative; and until such enumeration shall be made, the State of New Hampshire shall be entitled to chuse three, Massachusetts eight, Rhode Island and Providence Plantations one, Connecticut five, New York six, New Jersey four, Pennsylvania eight, Delaware one, Maryland six, Virginia ten, North Carolina five, South Carolina five and Georgia three. When vacancies happen in the Representation from any State, the Executive Authority thereof shall issue Writs of Election to fill such Vacancies. The House of Representatives shall chuse their Speaker and other Officers; and shall have the sole Power of Impeachment. Section 3. The Senate of the United States shall be composed of two Senators from each State, chosen by the Legislature thereof, [Superseded by Amendment XVII, section 1] for six Years; and each Senator shall have one Vote. Immediately after they shall be assembled in Consequence of the first Election, they shall be divided as equally as may be into three Classes. The Seats of the Senators of the first Class shall be vacated at the Expiration of the second Year, of the second Class at the Expiration of the fourth Year, and of the third Class at the Expiration of the sixth Year, so that one third may be chosen every second Year; and if Vacancies happen by Resignation, or otherwise, during the Recess of the Legislature of any State, the Executive thereof may make temporary Appointments until the next Meeting of the Legislature, which shall then fill such Vacancies. [Superseded by Amendment XVII, section 2] No person shall be a Senator who shall not have attained to the Age of thirty Years, and been nine Years a Citizen of the United States, and who shall not, when elected, be an Inhabitant of that State for which he shall be chosen. The Vice President of the United States shall be President of the Senate, but shall have no Vote, unless they be equally divided. The Senate shall choose their other Officers, and also a President pro tempore, in the absence of the Vice President, or when he shall exercise the Office of President of the United States. The Senate shall have the sole Power to try all Impeachments. When sitting for that Purpose, they shall be on Oath or Affirmation. When the President of the United States is tried, the Chief Justice shall preside: And no Person shall be convicted without the Concurrence of two thirds of the Members present. Judgment in Cases of Impeachment shall not extend further than to removal from Office, and disqualification to hold and enjoy any Office of honor, Trust or Profit under the United States: but the Party convicted shall nevertheless be liable and subject to Indictment, Trial, Judgment and Punishment, according to Law. Section 4. The Times, Places and Manner of holding Elections for Senators and Representatives, shall be prescribed in each State by the Legislature thereof; but the Congress may at any time by Law make or alter such Regulations, except as to the Place of Chusing Senators. The Congress shall assemble at least once in every Year, and such Meeting shall be on the first Monday in December, [Superseded by Amendment XX, section 2] unless they shall by Law appoint a different Day. Section 5. Each House shall be the Judge of the Elections, Returns and Qualifications of its own Members, and a Majority of each shall constitute a Quorum to do Business; but a smaller number may adjourn from day to day, and may be authorized to compel the Attendance of absent Members, in such Manner, and under such Penalties as each House may provide. Each House may determine the Rules of its Proceedings, punish its Members for disorderly Behavior, and, with the Concurrence of two-thirds, expel a Member. Each House shall keep a Journal of its Proceedings, and from time to time publish the same, excepting such Parts as may in their Judgment require Secrecy; and the Yeas and Nays of the Members of either House on any question shall, at the Desire of one fifth of those Present, be entered on the Journal. Neither House, during the Session of Congress, shall, without the Consent of the other, adjourn for more than three days, nor to any other Place than that in which the two Houses shall be sitting. Section 6. The Senators and Representatives shall receive a Compensation for their Services, to be ascertained by Law, and paid out of the Treasury of the United States. [Modified by Amendment XXVII] They shall in all Cases, except Treason, Felony and Breach of the Peace, be privileged from Arrest during their Attendance at the Session of their respective Houses, and in going to and returning from the same; and for any Speech or Debate in either House, they shall not be questioned in any other Place. No Senator or Representative shall, during the Time for which he was elected, be appointed to any civil Office under the Authority of the United States which shall have been created, or the Emoluments whereof shall have been increased during such time; and no Person holding any Office under the United States, shall be a Member of either House during his Continuance in Office. Section 7. All bills for raising Revenue shall originate in the House of Representatives; but the Senate may propose or concur with Amendments as on other Bills. Every Bill which shall have passed the House of Representatives and the Senate, shall, before it become a Law, be presented to the President of the United States; If he approve he shall sign it, but if not he shall return it, with his Objections to that House in which it shall have originated, who shall enter the Objections at large on their Journal, and proceed to reconsider it. If after such Reconsideration two thirds of that House shall agree to pass the Bill, it shall be sent, together with the Objections, to the other House, by which it shall likewise be reconsidered, and if approved by two thirds of that House, it shall become a Law. But in all such Cases the Votes of both Houses shall be determined by Yeas and Nays, and the Names of the Persons voting for and against the Bill shall be entered on the Journal of each House respectively. If any Bill shall not be returned by the President within ten Days (Sundays excepted) after it shall have been presented to him, the Same shall be a Law, in like Manner as if he had signed it, unless the Congress by their Adjournment prevent its Return, in which Case it shall not be a Law. Every Order, Resolution, or Vote to which the Concurrence of the Senate and House of Representatives may be necessary (except on a question of Adjournment) shall be presented to the President of the United States; and before the Same shall take Effect, shall be approved by him, or being disapproved by him, shall be repassed by two thirds of the Senate and House of Representatives, according to the Rules and Limitations prescribed in the Case of a Bill. Section 8. The Congress shall have Power To lay and collect Taxes, Duties, Imposts and Excises, to pay the Debts and provide for the common Defense and general Welfare of the United States; but all Duties, Imposts and Excises shall be uniform throughout the United States; To borrow money on the credit of the United States; To regulate Commerce with foreign Nations, and among the several States, and with the Indian Tribes; To establish an uniform Rule of Naturalization, and uniform Laws on the subject of Bankruptcies throughout the United States; To coin Money, regulate the Value thereof, and of foreign Coin, and fix the Standard of Weights and Measures; To provide for the Punishment of counterfeiting the Securities and current Coin of the United States; To establish Post Offices and Post Roads; To promote the Progress of Science and useful Arts, by securing for limited Times to Authors and Inventors the exclusive Right to their respective Writings and Discoveries; To constitute Tribunals inferior to the supreme Court; To define and punish Piracies and Felonies committed on the high Seas, and Offenses against the Law of Nations; To declare War, grant Letters of Marque and Reprisal, and make Rules concerning Captures on Land and Water; To raise and support Armies, but no Appropriation of Money to that Use shall be for a longer Term than two Years; To provide and maintain a Navy; To make Rules for the Government and Regulation of the land and naval Forces; To provide for calling forth the Militia to execute the Laws of the Union, suppress Insurrections and repel Invasions; To provide for organizing, arming, and disciplining the Militia, and for governing such Part of them as may be employed in the Service of the United States, reserving to the States respectively, the Appointment of the Officers, and the Authority of training the Militia according to the discipline prescribed by Congress; To exercise exclusive Legislation in all Cases whatsoever, over such District (not exceeding ten Miles square) as may, by Cession of particular States, and the acceptance of Congress, become the Seat of the Government of the United States, and to exercise like Authority over all Places purchased by the Consent of the Legislature of the State in which the Same shall be, for the Erection of Forts, Magazines, Arsenals, dock-Yards, and other needful Buildings; And To make all Laws which shall be necessary and proper for carrying into Execution the foregoing Powers, and all other Powers vested by this Constitution in the Government of the United States, or in any Department or Officer thereof. Section 9. The Migration or Importation of such Persons as any of the States now existing shall think proper to admit, shall not be prohibited by the Congress prior to the Year one thousand eight hundred and eight, but a tax or duty may be imposed on such Importation, not exceeding ten dollars for each Person. The privilege of the Writ of Habeas Corpus shall not be suspended, unless when in Cases of Rebellion or Invasion the public Safety may require it. No Bill of Attainder or ex post facto Law shall be passed. No capitation, or other direct, Tax shall be laid, unless in Proportion to the Census or Enumeration herein before directed to be taken. [Modified by Amendment XVI] No Tax or Duty shall be laid on Articles exported from any State. No Preference shall be given by any Regulation of Commerce or Revenue to the Ports of one State over those of another: nor shall Vessels bound to, or from, one State, be obliged to enter, clear, or pay Duties in another. No Money shall be drawn from the Treasury, but in Consequence of Appropriations made by Law; and a regular Statement and Account of the Receipts and Expenditures of all public Money shall be published from time to time. No Title of Nobility shall be granted by the United States: And no Person holding any Office of Profit or Trust under them, shall, without the Consent of the Congress, accept of any present, Emolument, Office, or Title, of any kind whatever, from any King, Prince or foreign State. Section 10. No State shall enter into any Treaty, Alliance, or Confederation; grant Letters of Marque and Reprisal; coin Money; emit Bills of Credit; make any Thing but gold and silver Coin a Tender in Payment of Debts; pass any Bill of Attainder, ex post facto Law, or Law impairing the Obligation of Contracts, or grant any Title of Nobility. No State shall, without the Consent of the Congress, lay any Imposts or Duties on Imports or Exports, except what may be absolutely necessary for executing it's inspection Laws: and the net Produce of all Duties and Imposts, laid by any State on Imports or Exports, shall be for the Use of the Treasury of the United States; and all such Laws shall be subject to the Revision and Controul of the Congress. No State shall, without the Consent of Congress, lay any duty of Tonnage, keep Troops, or Ships of War in time of Peace, enter into any Agreement or Compact with another State, or with a foreign Power, or engage in War, unless actually invaded, or in such imminent Danger as will not admit of delay. Article II. Section 1. The executive Power shall be vested in a President of the United States of America. He shall hold his Office during the Term of four Years, and, together with the Vice-President chosen for the same Term, be elected, as follows: Each State shall appoint, in such Manner as the Legislature thereof may direct, a Number of Electors, equal to the whole Number of Senators and Representatives to which the State may be entitled in the Congress: but no Senator or Representative, or Person holding an Office of Trust or Profit under the United States, shall be appointed an Elector. The Electors shall meet in their respective States, and vote by Ballot for two persons, of whom one at least shall not lie an Inhabitant of the same State with themselves. And they shall make a List of all the Persons voted for, and of the Number of Votes for each; which List they shall sign and certify, and transmit sealed to the Seat of the Government of the United States, directed to the President of the Senate. The President of the Senate shall, in the Presence of the Senate and House of Representatives, open all the Certificates, and the Votes shall then be counted. The Person having the greatest Number of Votes shall be the President, if such Number be a Majority of the whole Number of Electors appointed; and if there be more than one who have such Majority, and have an equal Number of Votes, then the House of Representatives shall immediately chuse by Ballot one of them for President; and if no Person have a Majority, then from the five highest on the List the said House shall in like Manner chuse the President. But in chusing the President, the Votes shall be taken by States, the Representation from each State having one Vote; a quorum for this Purpose shall consist of a Member or Members from two-thirds of the States, and a Majority of all the States shall be necessary to a Choice. In every Case, after the Choice of the President, the Person having the greatest Number of Votes of the Electors shall be the Vice President. But if there should remain two or more who have equal Votes, the Senate shall chuse from them by Ballot the Vice-President. The Congress may determine the Time of chusing the Electors, and the Day on which they shall give their Votes; which Day shall be the same throughout the United States. No person except a natural born Citizen, or a Citizen of the United States, at the time of the Adoption of this Constitution, shall be eligible to the Office of President; neither shall any Person be eligible to that Office who shall not have attained to the Age of thirty-five Years, and been fourteen Years a Resident within the United States. In Case of the Removal of the President from Office, or of his Death, Resignation, or Inability to discharge the Powers and Duties of the said Office, the same shall devolve on the Vice President, and the Congress may by Law provide for the Case of Removal, Death, Resignation or Inability, both of the President and Vice President, declaring what Officer shall then act as President, and such Officer shall act accordingly, until the Disability be removed, or a President shall be elected. The President shall, at stated Times, receive for his Services, a Compensation, which shall neither be increased nor diminished during the Period for which he shall have been elected, and he shall not receive within that Period any other Emolument from the United States, or any of them. Before he enter on the Execution of his Office, he shall take the following Oath or Affirmation: "I do solemnly swear (or affirm) that I will faithfully execute the Office of President of the United States, and will to the best of my Ability, preserve, protect and defend the Constitution of the United States." Section 2. The President shall be Commander in Chief of the Army and Navy of the United States, and of the Militia of the several States, when called into the actual Service of the United States; he may require the Opinion, in writing, of the principal Officer in each of the executive Departments, upon any subject relating to the Duties of their respective Offices, and he shall have Power to Grant Reprieves and Pardons for Offenses against the United States, except in Cases of Impeachment. He shall have Power, by and with the Advice and Consent of the Senate, to make Treaties, provided two thirds of the Senators present concur; and he shall nominate, and by and with the Advice and Consent of the Senate, shall appoint Ambassadors, other public Ministers and Consuls, Judges of the supreme Court, and all other Officers of the United States, whose Appointments are not herein otherwise provided for, and which shall be established by Law: but the Congress may by Law vest the Appointment of such inferior Officers, as they think proper, in the President alone, in the Courts of Law, or in the Heads of Departments. The President shall have Power to fill up all Vacancies that may happen during the Recess of the Senate, by granting Commissions which shall expire at the End of their next Session. Section 3. He shall from time to time give to the Congress Information of the State of the Union, and recommend to their Consideration such Measures as he shall judge necessary and expedient; he may, on extraordinary Occasions, convene both Houses, or either of them, and in Case of Disagreement between them, with Respect to the Time of Adjournment, he may adjourn them to such Time as he shall think proper; he shall receive Ambassadors and other public Ministers; he shall take Care that the Laws be faithfully executed, and shall Commission all the Officers of the United States. Section 4. The President, Vice President and all civil Officers of the United States, shall be removed from Office on Impeachment for, and Conviction of, Treason, Bribery, or other high Crimes and Misdemeanors. Article III. Section 1. The judicial Power of the United States, shall be vested in one supreme Court, and in such inferior Courts as the Congress may from time to time ordain and establish. The Judges, both of the supreme and inferior Courts, shall hold their Offices during good Behavior, and shall, at stated Times, receive for their Services a Compensation which shall not be diminished during their Continuance in Office. Section 2. The judicial Power shall extend to all Cases, in Law and Equity, arising under this Constitution, the Laws of the United States, and Treaties made, or which shall be made, under their Authority; to all Cases affecting Ambassadors, other public Ministers and Consuls; to all Cases of admiralty and maritime Jurisdiction; to Controversies to which the United States shall be a Party; to Controversies between two or more States; between a State and Citizens of another State; between Citizens of different States; between Citizens of the same State claiming Lands under Grants of different States, and between a State, or the Citizens thereof, and foreign States, Citizens or Subjects. In all Cases affecting Ambassadors, other public Ministers and Consuls, and those in which a State shall be Party, the supreme Court shall have original Jurisdiction. In all the other Cases before mentioned, the supreme Court shall have appellate Jurisdiction, both as to Law and Fact, with such Exceptions, and under such Regulations as the Congress shall make. The Trial of all Crimes, except in Cases of Impeachment, shall be by Jury; and such Trial shall be held in the State where the said Crimes shall have been committed; but when not committed within any State, the Trial shall be at such Place or Places as the Congress may by Law have directed. Section 3. Treason against the United States, shall consist only in levying War against them, or in adhering to their Enemies, giving them Aid and Comfort. No Person shall be convicted of Treason unless on the Testimony of two Witnesses to the same overt Act, or on Confession in open Court. The Congress shall have power to declare the Punishment of Treason, but no Attainder of Treason shall work Corruption of Blood, or Forfeiture except during the Life of the Person attainted. Article IV. Section 1. Full Faith and Credit shall be given in each State to the public Acts, Records, and judicial Proceedings of every other State. And the Congress may by general Laws prescribe the Manner in which such Acts, Records and Proceedings shall be proved, and the Effect thereof. Section 2. The Citizens of each State shall be entitled to all Privileges and Immunities of Citizens in the several States. A Person charged in any State with Treason, Felony, or other Crime, who shall flee from Justice, and be found in another State, shall on demand of the executive Authority of the State from which he fled, be delivered up, to be removed to the State having Jurisdiction of the Crime. No Person held to Service or Labour in one State, under the Laws thereof, escaping into another, shall, in Consequence of any Law or Regulation therein, be discharged from such Service or Labour, But shall be delivered up on Claim of the Party to whom such Service or Labour may be due. Section 3. New States may be admitted by the Congress into this Union; but no new States shall be formed or erected within the Jurisdiction of any other State; nor any State be formed by the Junction of two or more States, or parts of States, without the Consent of the Legislatures of the States concerned as well as of the Congress. The Congress shall have Power to dispose of and make all needful Rules and Regulations respecting the Territory or other Property belonging to the United States; and nothing in this Constitution shall be so construed as to Prejudice any Claims of the United States, or of any particular State. Section 4. The United States shall guarantee to every State in this Union a Republican Form of Government, and shall protect each of them against Invasion; and on Application of the Legislature, or of the Executive (when the Legislature cannot be convened) against domestic Violence. Article V. The Congress, whenever two thirds of both Houses shall deem it necessary, shall propose Amendments to this Constitution, or, on the Application of the Legislatures of two thirds of the several States, shall call a Convention for proposing Amendments, which, in either Case, shall be valid to all Intents and Purposes, as part of this Constitution, when ratified by the Legislatures of three fourths of the several States, or by Conventions in three fourths thereof, as the one or the other Mode of Ratification may be proposed by the Congress; Provided that no Amendment which may be made prior to the Year One thousand eight hundred and eight shall in any Manner affect the first and fourth Clauses in the Ninth Section of the first Article; and that no State, without its Consent, shall be deprived of its equal Suffrage in the Senate. Article VI. All Debts contracted and Engagements entered into, before the Adoption of this Constitution, shall be as valid against the United States under this Constitution, as under the Confederation. This Constitution, and the Laws of the United States which shall be made in Pursuance thereof; and all Treaties made, or which shall be made, under the Authority of the United States, shall be the supreme Law of the Land; and the Judges in every State shall be bound thereby, any Thing in the Constitution or Laws of any State to the Contrary notwithstanding. The Senators and Representatives before mentioned, and the Members of the several State Legislatures, and all executive and judicial Officers, both of the United States and of the several States, shall be bound by Oath or Affirmation, to support this Constitution; but no religious Test shall ever be required as a Qualification to any Office or public Trust under the United States. Article VII. The Ratification of the Conventions of nine States, shall be sufficient for the Establishment of this Constitution between the States so ratifying the Same. Closing endorsement. Done in Convention by the Unanimous Consent of the States present the Seventeenth Day of September in the Year of our Lord one thousand seven hundred and Eighty seven and of the Independence of the United States of America the Twelfth. In Witness whereof We have hereunto subscribed our Names. [Signatures Omitted] Amendment I. Congress shall make no law respecting an establishment of religion, or prohibiting the free exercise thereof; or abridging the freedom of speech, or of the press; or the right of the People peaceably to assemble, and to petition the Government for a redress of grievances. "Ratified December 15, 1791" Amendment II. A well regulated Militia, being necessary to the security of a free State, the right of the people to keep and bear Arms, shall not be infringed. "Ratified December 15, 1791" Amendment III. No Soldier shall, in time of peace be quartered in any house, without the consent of the Owner, nor in time of war, but in a manner to be prescribed by law. "Ratified December 15, 1791" Amendment IV. The right of the people to be secure in their persons, houses, papers, and effects, against unreasonable searches and seizures, shall not be violated, and no Warrants shall issue, but upon probable cause, supported by Oath or affirmation, and particularly describing the place to be searched, and the persons or things to be seized. "Ratified December 15, 1791" Amendment V. No person shall be held to answer for a capital, or otherwise infamous crime, unless on a presentment or indictment of a Grand Jury, except in cases arising in the land or naval forces, or in the Militia, when in actual service in time of War or public danger; nor shall any person be subject for the same offense to be twice put in jeopardy of life or limb; nor shall be compelled in any criminal case to be a witness against himself, nor be deprived of life, liberty, or property, without due process of law; nor shall private property be taken for public use, without just compensation. "Ratified December 15, 1791" Amendment VI. In all criminal prosecutions, the accused shall enjoy the right to a speedy and public trial, by an impartial jury of the State and district wherein the crime shall have been committed, which district shall have been previously ascertained by law, and to be informed of the nature and cause of the accusation; to be confronted with the witnesses against him; to have compulsory process for obtaining witnesses in his favor, and to have the Assistance of Counsel for his defence. "Ratified December 15, 1791" Amendment VII. In Suits at common law, where the value in controversy shall exceed twenty dollars, the right of trial by jury shall be preserved, and no fact tried by a jury, shall be otherwise re-examined in any Court of the United States, than according to the rules of the common law. "Ratified December 15, 1791" Amendment VIII. Excessive bail shall not be required, nor excessive fines imposed, nor cruel and unusual punishments inflicted. "Ratified December 15, 1791" Amendment IX. The enumeration in the Constitution, of certain rights, shall not be construed to deny or disparage others retained by the people. "Ratified December 15, 1791" Amendment X. The powers not delegated to the United States by the Constitution, nor prohibited by it to the States, are reserved to the States respectively, or to the people. "Ratified December 15, 1791" Amendment XI. The Judicial power of the United States shall not be construed to extend to any suit in law or equity, commenced or prosecuted against one of the United States by Citizens of another State, or by Citizens or Subjects of any Foreign State. "Ratified February 7, 1795" Amendment XII. The Electors shall meet in their respective states, and vote by ballot for President and Vice-President, one of whom, at least, shall not be an inhabitant of the same state with themselves; they shall name in their ballots the person voted for as President, and in distinct ballots the person voted for as Vice-President, and they shall make distinct lists of all persons voted for as President, and of all persons voted for as Vice-President and of the number of votes for each, which lists they shall sign and certify, and transmit sealed to the seat of the government of the United States, directed to the President of the Senate; The President of the Senate shall, in the presence of the Senate and House of Representatives, open all the certificates and the votes shall then be counted; The person having the greatest Number of votes for President, shall be the President, if such number be a majority of the whole number of Electors appointed; and if no person have such majority, then from the persons having the highest numbers not exceeding three on the list of those voted for as President, the House of Representatives shall choose immediately, by ballot, the President. But in choosing the President, the votes shall be taken by states, the representation from each state having one vote; a quorum for this purpose shall consist of a member or members from two-thirds of the states, and a majority of all the states shall be necessary to a choice. And if the House of Representatives shall not choose a President whenever the right of choice shall devolve upon them, before the fourth day of March next following, then the Vice-President shall act as President, as in the case of the death or other constitutional disability of the President. The person having the greatest number of votes as Vice-President, shall be the Vice-President, if such number be a majority of the whole number of Electors appointed, and if no person have a majority, then from the two highest numbers on the list, the Senate shall choose the Vice-President; a quorum for the purpose shall consist of two-thirds of the whole number of Senators, and a majority of the whole number shall be necessary to a choice. But no person constitutionally ineligible to the office of President shall be eligible to that of Vice-President of the United States. "Ratified June 15, 1804" Amendment XIII. Section 1. Neither slavery nor involuntary servitude, except as a punishment for crime whereof the party shall have been duly convicted, shall exist within the United States &amp; all of its territories, or any place subject to their jurisdiction. Section 2. Congress shall have power to enforce this article by appropriate legislation. "Ratified December 6, 1865" Amendment XIV. Section 1. All persons born or naturalized in the United States, and subject to the jurisdiction thereof, are citizens of the United States and of the State wherein they reside. No State shall make or enforce any law which shall abridge the privileges or immunities of citizens of the United States; nor shall any State deprive any person of life, liberty, or property, without due process of law; nor deny to any person within its jurisdiction the equal protection of the laws. Section 2. Representatives shall be apportioned among the several States according to their respective numbers, counting the whole number of persons in each State, excluding Indians not taxed. But when the right to vote at any election for the choice of electors for President and Vice-President of the United States, Representatives in Congress, the Executive and Judicial officers of a State, or the members of the Legislature thereof, is denied to any of the male inhabitants of such State, being twenty-one years of age, and citizens of the United States, or in any way abridged, except for participation in rebellion, or other crime, the basis of representation therein shall be reduced in the proportion which the number of such male citizens shall bear to the whole number of male citizens twenty-one years of age in such State. Section 3. No person shall be a Senator or Representative in Congress, or elector of President and Vice-President, or hold any office, civil or military, under the United States, or under any State, who, having previously taken an oath, as a member of Congress, or as an officer of the United States, or as a member of any State legislature, or as an executive or judicial officer of any State, to support the Constitution of the United States, shall have engaged in insurrection or rebellion against the same, or given aid or comfort to the enemies thereof. But Congress may by a vote of two-thirds of each House, remove such disability. Section 4. The validity of the public debt of the United States, authorized by law, including debts incurred for payment of pensions and bounties for services in suppressing insurrection or rebellion, shall not be questioned. But neither the United States nor any State shall assume or pay any debt or obligation incurred in aid of insurrection or rebellion against the United States, or any claim for the loss or emancipation of any slave; but all such debts, obligations and claims shall be held illegal and void. Section 5. The Congress shall have power to enforce, by appropriate legislation, the provisions of this article. "Ratified July 9, 1868" Amendment XV. Section 1. The right of citizens of the United States to vote shall not be denied or abridged by the United States or by any State on account of race, color, or previous condition of servitude. Section 2. The Congress shall have power to enforce this article by appropriate legislation. "Ratified February 3, 1870" Amendment XVI. The Congress shall have power to lay and collect taxes on incomes, from whatever source derived, without apportionment among the several States, and without regard to any census or enumeration. "Ratified February 3, 1913" Amendment XVII. The Senate of the United States shall be composed of two Senators from each State, elected by the people thereof, for six years; and each Senator shall have one vote. The electors in each State shall have the qualifications requisite for electors of the most numerous branch of the State legislatures. When vacancies happen in the representation of any State in the Senate, the executive authority of such State shall issue writs of election to fill such vacancies: Provided, That the legislature of any State may empower the executive thereof to make temporary appointments until the people fill the vacancies by election as the legislature may direct. This amendment shall not be so construed as to affect the election or term of any Senator chosen before it becomes valid as part of the Constitution. "Ratified April 8, 1913" Amendment XVIII. Section 1. After one year from the ratification of this article the manufacture, sale, or transportation of intoxicating liquors within, the importation thereof into, or the exportation thereof from the United States and all territory subject to the jurisdiction thereof for beverage purposes is hereby prohibited. Section 2. The Congress and the several States shall have concurrent power to enforce this article by appropriate legislation. Section 3. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of the several States, as provided in the Constitution, within seven years from the date of the submission hereof to the States by the Congress. "Ratified January 16, 1919. Repealed by Amendment XXI December 5, 1933" Amendment XIX. The right of citizens of the United States to vote shall not be denied or abridged by the United States or by any State on account of sex. Congress shall have power to enforce this article by appropriate legislation. "Ratified August 18, 1920" Amendment XX. Section 1. The terms of the President and Vice President shall end at noon on the 20th day of January, and the terms of Senators and Representatives at noon on the 3rd day of January, of the years in which such terms would have ended if this article had not been ratified; and the terms of their successors shall then begin. Section 2. The Congress shall assemble at least once in every year, and such meeting shall begin at noon on the 3rd day of January, unless they shall by law appoint a different day. Section 3. If, at the time fixed for the beginning of the term of the President, the President elect shall have died, the Vice President elect shall become President. If a President shall not have been chosen before the time fixed for the beginning of his term, or if the President elect shall have failed to qualify, then the Vice President elect shall act as President until a President shall have qualified; and the Congress may by law provide for the case wherein neither a President elect nor a Vice President elect shall have qualified, declaring who shall then act as President, or the manner in which one who is to act shall be selected, and such person shall act accordingly until a President or Vice President shall have qualified. Section 4. The Congress may by law provide for the case of the death of any of the persons from whom the House of Representatives may choose a President whenever the right of choice shall have devolved upon them, and for the case of the death of any of the persons from whom the Senate may choose a Vice President whenever the right of choice shall have devolved upon them. Section 5. Sections 1 and 2 shall take effect on the 15th day of October following the ratification of this article. Section 6. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of three-fourths of the several States within seven years from the date of its submission. "Ratified January 23, 1933" Amendment XXI. Section 1. The eighteenth article of amendment to the Constitution of the United States is hereby repealed. Section 2. The transportation or importation into any State, Territory, or possession of the United States for delivery or use therein of intoxicating liquors, in violation of the laws thereof, is hereby prohibited. Section 3. The article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by conventions in the several States, as provided in the Constitution, within seven years from the date of the submission hereof to the States by the Congress. "Ratified December 5, 1933" Amendment XXII. Section 1. No person shall be elected to the office of the President more than twice, and no person who has held the office of President, or acted as President, for more than two years of a term to which some other person was elected President shall be elected to the office of the President more than once. But this Article shall not apply to any person holding the office of President, when this Article was proposed by the Congress, and shall not prevent any person who may be holding the office of President, or acting as President, during the term within which this Article becomes operative from holding the office of President or acting as President during the remainder of such term. Section 2. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of three-fourths of the several States within seven years from the date of its submission to the States by the Congress. "Ratified February 27, 1951" Amendment XXIII. Section 1. The District constituting the seat of Government of the United States shall appoint in such manner as the Congress may direct: A number of electors of President and Vice President equal to the whole number of Senators and Representatives in Congress to which the District would be entitled if it were a State, but in no event more than the least populous State; they shall be in addition to those appointed by the States, but they shall be considered, for the purposes of the election of President and Vice President, to be electors appointed by a State; and they shall meet in the District and perform such duties as provided by the twelfth article of amendment. Section 2. The Congress shall have power to enforce this article by appropriate legislation. "Ratified March 29, 1961" Amendment XXIV. Section 1. The right of citizens of the United States to vote in any primary or other election for President or Vice President, for electors for President or Vice President, or for Senator or Representative in Congress, shall not be denied or abridged by the United States or any State by reason of failure to pay any poll tax or other tax. Section 2. The Congress shall have power to enforce this article by appropriate legislation. "Ratified January 23, 1964" Amendment XXV. Section 1. In case of the removal of the President from office or of his death or resignation, the Vice President shall become President. Section 2. Whenever there is a vacancy in the office of the Vice President, the President shall nominate a Vice President who shall take office upon confirmation by a majority vote of both Houses of Congress. Section 3. Whenever the President transmits to the President pro tempore of the Senate and the Speaker of the House of Representatives his written declaration that he is unable to discharge the powers and duties of his office, and until he transmits to them a written declaration to the contrary, such powers and duties shall be discharged by the Vice President as Acting President. Section 4. Whenever the Vice President and a majority of either the principal officers of the executive departments or of such other body as Congress may by law provide, transmit to the President pro tempore of the Senate and the Speaker of the House of Representatives their written declaration that the President is unable to discharge the powers and duties of his office, the Vice President shall immediately assume the powers and duties of the office as Acting President. Thereafter, when the President transmits to the President pro tempore of the Senate and the Speaker of the House of Representatives his written declaration that no inability exists, he shall resume the powers and duties of his office unless the Vice President and a majority of either the principal officers of the executive department or of such other body as Congress may by law provide, transmit within four days to the President pro tempore of the Senate and the Speaker of the House of Representatives their written declaration that the President is unable to discharge the powers and duties of his office. Thereupon Congress shall decide the issue, assembling within forty eight hours for that purpose if not in session. If the Congress, within twenty one days after receipt of the latter written declaration, or, if Congress is not in session, within twenty one days after Congress is required to assemble, determines by two thirds vote of both Houses that the President is unable to discharge the powers and duties of his office, the Vice President shall continue to discharge the same as Acting President; otherwise, the President shall resume the powers and duties of his office. "Ratified February 10, 1967" Amendment XXVI. Section 1. The right of citizens of the United States, who are eighteen years of age or older, to vote shall not be denied or abridged by the United States or by any State on account of age. Section 2. The Congress shall have power to enforce this article by appropriate legislation. "Ratified July 1, 1971" Amendment XXVII. No law, varying the compensation for the services of the Senators and Representatives, shall take effect, until an election of Representatives shall have intervened. "Ratified May 7, 1992" 

In Congress, July 4, 1776 The unanimous Declaration of the thirteen united States of America When in the Course of human events, it becomes necessary for one people to dissolve the political bands which have connected them with another, and to assume among the powers of the earth, the separate and equal station to which the Laws of Nature and of Nature's God entitle them, a decent respect to the opinions of mankind requires that they should declare the causes which impel them to the separation. We hold these truths to be self-evident, that all men are created equal, that they are endowed by their Creator with certain unalienable Rights, that among these are Life, Liberty and the pursuit of Happiness. That to secure these rights, Governments are instituted among Men, deriving their just Powers from the consent of the governed, -- That whenever any Form of Government becomes destructive of these ends, it is the Right of the People to alter or to abolish it, and to institute new Government, laying its foundation on such principles and organizing its powers in such form, as to them shall seem most likely to effect their Safety and Happiness. Prudence, indeed, will dictate that Governments long established should not be changed for light and transient causes; and accordingly all experience hath shewn, that mankind are more disposed to suffer, while evils are sufferable, than to right themselves by abolishing the forms to which they are accustomed. But when a long train of abuses and usurpations, pursuing invariably the same Object evinces a design to reduce them under absolute Despotism, it is their right, it is their duty, to throw off such Government, and to provide new guards for their future security -- Such has been the patient sufferance of these Colonies; and such is now the necessity which constrains them to alter their former Systems of Government. -- The history of the present King of Great Britain is a history of repeated injuries and usurpations, all having in direct object the establishment of an absolute Tyranny over these States. To prove this, let facts be submitted to a candid world. He has refused his Assent to Laws, the most wholesome and necessary for the public good. He has forbidden his Governors to pass Laws of immediate and pressing importance, unless suspended in their operation till his Assent should be obtained; and when so suspended, he has utterly neglected to attend to them. He has refused to pass other Laws for the accommodation of large districts of people, unless those people would relinquish the right of Representation in the Legislature, a right inestimable to them and formidable to tyrants only. He has called together legislative bodies at places unusual, uncomfortable, and distant from the depository of their Public Records, for the sole purpose of fatiguing them into compliance with his measures. He has dissolved Representative Houses repeatedly, for opposing with manly firmness his invasions on the rights of the people. He has refused for a long time, after such dissolutions, to cause others to be elected; whereby the Legislative Powers, incapable of Annihilation, have returned to the People at large for their exercise; the State remaining in the mean time exposed to all the dangers of invasion from without, and convulsions within. He has endeavoured to prevent the population of these States; for that purpose obstructing the Laws for Naturalization of Foreigners; refusing to pass others to encourage their migrations hither, and raising the conditions of new Appropriations of Lands. He has obstructed the Administration of Justice, by refusing his Assent to Laws for establishing Judiciary Powers. He has made Judges dependent on his Will alone, for the tenure of their offices, and the amount and payment of their salaries. He has erected a multitude of New Offices, and sent hither swarms of Officers to harrass our People, and eat out their substance. He has kept among us, in times of peace, Standing Armies without the Consent of our legislatures. He has affected to render the Military independent of and superior to the Civil Power. He has combined with others to subject us to a jurisdiction foreign to our constitution, and unacknowledged by our laws; giving his Assent to their Acts of pretended Legislation: For Quartering large bodies of armed troops among us: For protecting them, by a mock Trial, from Punishment for any Murders which they should commit on the Inhabitants of these States: For cutting off our Trade with all parts of the world: For imposing Taxes on us without our Consent: For depriving us in many cases, of the benefits of Trial by Jury: For transporting us beyond seas to be tried for pretended offences: For abolishing the free system of English Laws in a neighbouring Province, establishing therein an Arbitrary government, and enlarging its Boundaries so as to render it at once an example and fit instrument for introducing the same absolute rule into these Colonies: For taking away our Charters, abolishing our most valuable Laws, and altering fundamentally the forms of our Governments: For suspending our own Legislature, and declaring themselves invested with power to legislate for us in all cases whatsoever. He has abdicated Government here, by declaring us out of his Protection and waging War against us. He has plundered our seas, ravaged our Coasts, burnt our towns, and destroyed the lives of our people. He is at this time transporting large Armies of foreign Mercenaries to compleat the works of death, desolation and tyranny, already begun with circumstances of Cruelty and perfidy scarcely paralleled in the most barbarous ages, and totally unworthy the Head of a civilized nation. He has constrained our fellow Citizens taken Captive on the high Seas to bear Arms against their Country, to become the executioners of their friends and Brethren, or to fall themselves by their Hands. He has excited domestic insurrections amongst us, and has endeavoured to bring on the inhabitants of our frontiers, the merciless Indian Savages, whose known rule of warfare, is an undistinguished destruction of all ages, sexes and conditions. In every stage of these Oppressions we have Petitioned for Redress in the most humble terms: Our repeated Petitions have been answered only by repeated injury. A Prince, whose character is thus marked by every act which may define a Tyrant, is unfit to be the ruler of a free people. Nor have we been wanting in attention to our British brethren. We have warned them from time to time of attempts by their legislature to extend an unwarrantable jurisdiction over us. We have reminded them of the circumstances of our emigration and settlement here. We have appealed to their native justice and magnanimity, and we have conjured them by the ties of our common kindred to disavow these usurpations, which, would inevitably interrupt our connections and correspondence. They too have been deaf to the voice of justice and of consanguinity. We must, therefore, acquiesce in the necessity, which denounces our Separation, and hold them, as we hold the rest of mankind, Enemies in War, in Peace Friends. We, therefore, the Representatives of the united States of America, in General Congress, Assembled, appealing to the Supreme Judge of the world for the rectitude of our intentions, do, in the Name, and by Authority of the good People of these Colonies, solemnly publish and declare, That these United Colonies are, and of Right ought to be Free and Independent States; that they are absolved from all Allegiance to the British Crown, and that all political connection between them and the State of Great Britain, is and ought to be totally dissolved; and that as Free and Independent States, they have full Power to levy War, conclude Peace, contract Alliances, establish Commerce, and to do all other Acts and Things which Independent States may of right do. And for the support of this Declaration, with a firm reliance on the protection of Divine Providence, we mutually pledge to each other our Lives, our Fortunes and our sacred Honor. [Signatures Omitted] 

Preamble. To all to whom these Presents shall come, we the undersigned Delegates of the States affixed to our Names send greeting. Articles of Confederation and perpetual Union between the States of New Hampshire, Massachusetts bay, Rhode Island and Providence Plantations, Connecticut, New York, New Jersey, Pennsylvania, Delaware, Maryland, Virginia, North Carolina, South Carolina and Georgia. Article I. The Stile of this Confederacy shall be "The United States of America." Article II. Each state retains its sovereignty, freedom, and independence, and every power, jurisdiction, and right, which is not by this Confederation expressly delegated to the United States, in Congress assembled. Article III. The said States hereby severally enter into a firm league of friendship with each other, for their common defense, the security of their liberties, and their mutual and general welfare, binding themselves to assist each other, against all force offered to, or attacks made upon them, or any of them, on account of religion, sovereignty, trade, or any other pretense whatever. Article IV. The better to secure and perpetuate mutual friendship and intercourse among the people of the different States in this Union, the free inhabitants of each of these States, paupers, vagabonds, and fugitives from justice excepted, shall be entitled to all privileges and immunities of free citizens in the several States; and the people of each State shall free ingress and regress to and from any other State, and shall enjoy therein all the privileges of trade and commerce, subject to the same duties, impositions, and restrictions as the inhabitants thereof respectively, provided that such restrictions shall not extend so far as to prevent the removal of property imported into any State, to any other State, of which the owner is an inhabitant; provided also that no imposition, duties or restriction shall be laid by any State, on the property of the United States, or either of them. If any person guilty of, or charged with, treason, felony, or other high misdemeanor in any State, shall flee from justice, and be found in any of the United States, he shall, upon demand of the Governor or executive power of the State from which he fled, be delivered up and removed to the State having jurisdiction of his offense. Full faith and credit shall be given in each of these States to the records, acts, and judicial proceedings of the courts and magistrates of every other State. Article V. For the most convenient management of the general interests of the United States, delegates shall be annually appointed in such manner as the legislatures of each State shall direct, to meet in Congress on the first Monday in November, in every year, with a power reserved to each State to recall its delegates, or any of them, at any time within the year, and to send others in their stead for the remainder of the year. No State shall be represented in Congress by less than two, nor more than seven members; and no person shall be capable of being a delegate for more than three years in any term of six years; nor shall any person, being a delegate, be capable of holding any office under the United States, for which he, or another for his benefit, receives any salary, fees or emolument of any kind. Each State shall maintain its own delegates in a meeting of the States, and while they act as members of the committee of the States. In determining questions in the United States in Congress assembled, each State shall have one vote. Freedom of speech and debate in Congress shall not be impeached or questioned in any court or place out of Congress, and the members of Congress shall be protected in their persons from arrests or imprisonments, during the time of their going to and from, and attendance on Congress, except for treason, felony, or breach of the peace. Article VI. No State, without the consent of the United States in Congress assembled, shall send any embassy to, or receive any embassy from, or enter into any conference, agreement, alliance or treaty with any King, Prince or State; nor shall any person holding any office of profit or trust under the United States, or any of them, accept any present, emolument, office or title of any kind whatever from any King, Prince or foreign State; nor shall the United States in Congress assembled, or any of them, grant any title of nobility. No two or more States shall enter into any treaty, confederation or alliance whatever between them, without the consent of the United States in Congress assembled, specifying accurately the purposes for which the same is to be entered into, and how long it shall continue. No State shall lay any imposts or duties, which may interfere with any stipulations in treaties, entered into by the United States in Congress assembled, with any King, Prince or State, in pursuance of any treaties already proposed by Congress, to the courts of France and Spain. No vessel of war shall be kept up in time of peace by any State, except such number only, as shall be deemed necessary by the United States in Congress assembled, for the defense of such State, or its trade; nor shall any body of forces be kept up by any State in time of peace, except such number only, as in the judgement of the United States in Congress assembled, shall be deemed requisite to garrison the forts necessary for the defense of such State; but every State shall always keep up a well-regulated and disciplined militia, sufficiently armed and accoutered, and shall provide and constantly have ready for use, in public stores, a due number of filed pieces and tents, and a proper quantity of arms, ammunition and camp equipage. No State shall engage in any war without the consent of the United States in Congress assembled, unless such State be actually invaded by enemies, or shall have received certain advice of a resolution being formed by some nation of Indians to invade such State, and the danger is so imminent as not to admit of a delay till the United States in Congress assembled can be consulted; nor shall any State grant commissions to any ships or vessels of war, nor letters of marque or reprisal, except it be after a declaration of war by the United States in Congress assembled, and then only against the Kingdom or State and the subjects thereof, against which war has been so declared, and under such regulations as shall be established by the United States in Congress assembled, unless such State be infested by pirates, in which case vessels of war may be fitted out for that occasion, and kept so long as the danger shall continue, or until the United States in Congress assembled shall determine otherwise. Article VII. When land forces are raised by any State for the common defense, all officers of or under the rank of colonel, shall be appointed by the legislature of each State respectively, by whom such forces shall be raised, or in such manner as such State shall direct, and all vacancies shall be filled up by the State which first made the appointment. Article VIII. All charges of war, and all other expenses that shall be incurred for the common defense or general welfare, and allowed by the United States in Congress assembled, shall be defrayed out of a common treasury, which shall be supplied by the several States in proportion to the value of all land within each State, granted or surveyed for any person, as such land and the buildings and improvements thereon shall be estimated according to such mode as the United States in Congress assembled, shall from time to time direct and appoint. The taxes for paying that proportion shall be laid and levied by the authority and direction of the legislatures of the several States within the time agreed upon by the United States in Congress assembled. Article IX. The United States in Congress assembled, shall have the sole and exclusive right and power of determining on peace and war, except in the cases mentioned in the sixth article -- of sending and receiving ambassadors -- entering into treaties and alliances, provided that no treaty of commerce shall be made whereby the legislative power of the respective States shall be restrained from imposing such imposts and duties on foreigners, as their own people are subjected to, or from prohibiting the exportation or importation of any species of goods or commodities whatsoever -- of establishing rules for deciding in all cases, what captures on land or water shall be legal, and in what manner prizes taken by land or naval forces in the service of the United States shall be divided or appropriated -- of granting letters of marque and reprisal in times of peace -- appointing courts for the trial of piracies and felonies committed on the high seas and establishing courts for receiving and determining finally appeals in all cases of captures, provided that no member of Congress shall be appointed a judge of any of the said courts. The United States in Congress assembled shall also be the last resort on appeal in all disputes and differences now subsisting or that hereafter may arise between two or more States concerning boundary, jurisdiction or any other causes whatever; which authority shall always be exercised in the manner following. Whenever the legislative or executive authority or lawful agent of any State in controversy with another shall present a petition to Congress stating the matter in question and praying for a hearing, notice thereof shall be given by order of Congress to the legislative or executive authority of the other State in controversy, and a day assigned for the appearance of the parties by their lawful agents, who shall then be directed to appoint by joint consent, commissioners or judges to constitute a court for hearing and determining the matter in question: but if they cannot agree, Congress shall name three persons out of each of the United States, and from the list of such persons each party shall alternately strike out one, the petitioners beginning, until the number shall be reduced to thirteen; and from that number not less than seven, nor more than nine names as Congress shall direct, shall in the presence of Congress be drawn out by lot, and the persons whose names shall be so drawn or any five of them, shall be commissioners or judges, to hear and finally determine the controversy, so always as a major part of the judges who shall hear the cause shall agree in the determination: and if either party shall neglect to attend at the day appointed, without showing reasons, which Congress shall judge sufficient, or being present shall refuse to strike, the Congress shall proceed to nominate three persons out of each State, and the secretary of Congress shall strike in behalf of such party absent or refusing; and the judgement and sentence of the court to be appointed, in the manner before prescribed, shall be final and conclusive; and if any of the parties shall refuse to submit to the authority of such court, or to appear or defend their claim or cause, the court shall nevertheless proceed to pronounce sentence, or judgement, which shall in like manner be final and decisive, the judgement or sentence and other proceedings being in either case transmitted to Congress, and lodged among the acts of Congress for the security of the parties concerned: provided that every commissioner, before he sits in judgement, shall take an oath to be administered by one of the judges of the supreme or superior court of the State, where the cause shall be tried, 'well and truly to hear and determine the matter in question, according to the best of his judgement, without favor, affection or hope of reward': provided also, that no State shall be deprived of territory for the benefit of the United States. All controversies concerning the private right of soil claimed under different grants of two or more States, whose jurisdictions as they may respect such lands, and the States which passed such grants are adjusted, the said grants or either of them being at the same time claimed to have originated antecedent to such settlement of jurisdiction, shall on the petition of either party to the Congress of the United States, be finally determined as near as may be in the same manner as is before prescribed for deciding disputes respecting territorial jurisdiction between different States. The United States in Congress assembled shall also have the sole and exclusive right and power of regulating the alloy and value of coin struck by their own authority, or by that of the respective States -- fixing the standards of weights and measures throughout the United States -- regulating the trade and managing all affairs with the Indians, not members of any of the States, provided that the legislative right of any State within its own limits be not infringed or violated -- establishing or regulating post offices from one State to another, throughout all the United States, and exacting such postage on the papers passing through the same as may be requisite to defray the expenses of the said office -- appointing all officers of the land forces, in the service of the United States, excepting regimental officers -- appointing all the officers of the naval forces, and commissioning all officers whatever in the service of the United States -- making rules for the government and regulation of the said land and naval forces, and directing their operations. The United States in Congress assembled shall have authority to appoint a committee, to sit in the recess of Congress, to be denominated 'A Committee of the States', and to consist of one delegate from each State; and to appoint such other committees and civil officers as may be necessary for managing the general affairs of the United States under their direction -- to appoint one of their members to preside, provided that no person be allowed to serve in the office of president more than one year in any term of three years; to ascertain the necessary sums of money to be raised for the service of the United States, and to appropriate and apply the same for defraying the public expenses -- to borrow money, or emit bills on the credit of the United States, transmitting every half-year to the respective States an account of the sums of money so borrowed or emitted -- to build and equip a navy -- to agree upon the number of land forces, and to make requisitions from each State for its quota, in proportion to the number of white inhabitants in such State; which requisition shall be binding, and thereupon the legislature of each State shall appoint the regimental officers, raise the men and cloath, arm and equip them in a solid- like manner, at the expense of the United States; and the officers and men so cloathed, armed and equipped shall march to the place appointed, and within the time agreed on by the United States in Congress assembled. But if the United States in Congress assembled shall, on consideration of circumstances judge proper that any State should not raise men, or should raise a smaller number of men than the quota thereof, such extra number shall be raised, officered, cloathed, armed and equipped in the same manner as the quota of each State, unless the legislature of such State shall judge that such extra number cannot be safely spread out in the same, in which case they shall raise, officer, cloath, arm and equip as many of such extra number as they judge can be safely spared. And the officers and men so cloathed, armed, and equipped, shall march to the place appointed, and within the time agreed on by the United States in Congress assembled. The United States in Congress assembled shall never engage in a war, nor grant letters of marque or reprisal in time of peace, nor enter into any treaties or alliances, nor coin money, nor regulate the value thereof, nor ascertain the sums and expenses necessary for the defense and welfare of the United States, or any of them, nor emit bills, nor borrow money on the credit of the United States, nor appropriate money, nor agree upon the number of vessels of war, to be built or purchased, or the number of land or sea forces to be raised, nor appoint a commander in chief of the army or navy, unless nine States assent to the same: nor shall a question on any other point, except for adjourning from day to day be determined, unless by the votes of the majority of the United States in Congress assembled. The Congress of the United States shall have power to adjourn to any time within the year, and to any place within the United States, so that no period of adjournment be for a longer duration than the space of six months, and shall publish the journal of their proceedings monthly, except such parts thereof relating to treaties, alliances or military operations, as in their judgement require secrecy; and the yeas and nays of the delegates of each State on any question shall be entered on the journal, when it is desired by any delegates of a State, or any of them, at his or their request shall be furnished with a transcript of the said journal, except such parts as are above excepted, to lay before the legislatures of the several States. Article X. The Committee of the States, or any nine of them, shall be authorized to execute, in the recess of Congress, such of the powers of Congress as the United States in Congress assembled, by the consent of the nine States, shall from time to time think expedient to vest them with; provided that no power be delegated to the said Committee, for the exercise of which, by the Articles of Confederation, the voice of nine States in the Congress of the United States assembled be requisite. Article XI. Canada acceding to this confederation, and adjoining in the measures of the United States, shall be admitted into, and entitled to all the advantages of this Union; but no other colony shall be admitted into the same, unless such admission be agreed to by nine States. Article XII. All bills of credit emitted, monies borrowed, and debts contracted by, or under the authority of Congress, before the assembling of the United States, in pursuance of the present confederation, shall be deemed and considered as a charge against the United States, for payment and satisfaction whereof the said United States, and the public faith are hereby solemnly pledged. Article XIII. Every State shall abide by the determination of the United States in Congress assembled, on all questions which by this confederation are submitted to them. And the Articles of this Confederation shall be inviolably observed by every State, and the Union shall be perpetual; nor shall any alteration at any time hereafter be made in any of them; unless such alteration be agreed to in a Congress of the United States, and be afterwards confirmed by the legislatures of every State. Conclusion. And Whereas it hath pleased the Great Governor of the World to incline the hearts of the legislatures we respectively represent in Congress, to approve of, and to authorize us to ratify the said Articles of Confederation and perpetual Union. Know Ye that we the undersigned delegates, by virtue of the power and authority to us given for that purpose, do by these presents, in the name and in behalf of our respective constituents, fully and entirely ratify and confirm each and every of the said Articles of Confederation and perpetual Union, and all and singular the matters and things therein contained: And we do further solemnly plight and engage the faith of our respective constituents, that they shall abide by the determinations of the United States in Congress assembled, on all questions, which by the said Confederation are submitted to them. And that the Articles thereof shall be inviolably observed by the States we respectively represent, and that the Union shall be perpetual. [Signatures Omitted] 

Introduction. Since the Bill of Rights, there have been seventeen Amendments to the Constitution. Eleventh Amendment. "The Judicial power of the United States shall not be construed to extend to any suit in law or equity, commenced or prosecuted against one of the United States by Citizens of another State, or by Citizens or Subjects of any Foreign State." The eleventh amendment was adopted in response to the Supreme Court case "Chisolm v. Georgia", (1793) in which the court ruled that the Constitution granted federal courts the power to hear cases brought against states from its own citizens or citizens of different state or citizens or subject of foreign countries. This established that states lack sovereign immunity from suits brought in federal courts. The 11th amendment however, superseded and overruled the supreme court ruling and established that the states have sovereign immunity from from suits brought against them from citizens of another states or citizens of foreign countries. Twelfth Amendment. The Twelfth Amendment fine-tuned the process for electing the President. (See Part III, Chapter 2 for details.) Thirteenth Amendment. The Thirteenth Amendment abolished and outlawed slavery everywhere in the United States and every place subject to their jurisdiction. It also abolished and outlawed involuntary servitude except to punish crime. Fourteenth Amendment. Amended in 1868, the Fourteenth Amendment served to extend much of the Bill of Rights to the states by requiring that "No State shall make or enforce any law which shall abridge the privileges or immunities of citizens of the United States." It was essentially a reaffirmation that all citizens are considered equal regardless of race. Secondly, the Amendment provided that a State cannot deprive any person of life, liberty or property without due process of law. Thirdly, the amendment required that every State provide equal protection to all of its citizens- this clause intended to prevent discrimination against African-Americans, although due to several supreme court rulings the clauses effects were almost rendered invalid. Fourthly, it removed the original Constitutional requirement that "other persons" (slaves) count as three-fifths of a person when determining the official state population and that all persons count toward a States population. Fifthly, it stated that a State's representation in Congress would be reduced if that State prohibited males over twenty-one from voting for a reason other than commission of a crime. Sixthly, the Amendment prohibited any person who participates in a rebellion such as the Civil War from serving in government unless the Congress formally agrees, by a two-thirds vote in each house, to exempt the person from this disqualification. Seventhly, the Fourteenth required that any debts incurred by the Union (the North) during the Civil War were to be held valid and as such had to be paid, but that any and all debts incurred by the Confederacy were illegal and void and as such could not and would not be paid. It also stated that any person who lost slaves during the civil war could not sue the Government because of the loss thereof. Fifteenth Amendment. The Fifteenth Amendment took another barrier away from minorities right to vote by barring the creation of laws that barred someone from voting based on their color, race, or having been former slaves. Sixteenth Amendment. The Sixteenth Amendment made Income Taxes constitutional. It was passed to resolve disputes regarding the matter. At one time, the Supreme Court ruled that the Tax was constitutional, while ruling at another time that it was not. All doubts were removed by the Sixteenth. Seventeenth Amendment. The Seventeenth provided that the people of a state, and not the legislature, would be the electors of Senators. However, it retained the power of legislatures to temporarily fill vacancies until an election can be held. Eighteenth Amendment. The Eighteenth Amendment created Prohibition, banning intoxicating liquors in the United States. It allowed both the federal and the state governments to concurrently legislate on the matter. Nineteenth Amendment. The Nineteenth granted women the right to vote on an equal basis with men. Twentieth Amendment. The Twentieth Amendment was known as the "lame duck" amendment. According to the original constitution and the twelfth amendment, a newly elected president did not take office until March 4. In the early days of the Republic, elections were held on different days in different jurisdictions, and travel and communication were slow. By the 1930s, however, the election was uniformly held in early November, and the winner of a presidential election was usually apparent by the next morning. This meant that a newly elected president had to wait four months take take office, leavin the outgoing president a "lame duck." This amendment moved up the date for a new president to take office from March 4 to January 20. The first day for the newly elected Congress was moved up to January 3, allowing Congress time to select the president or vice president should the need arise. The amendment also required that Congress meet each year on January 3, unless a different day was chosen. The amendment also provided for several remote circumstances involving deaths or resignations of presidents-elect and vice presidents-elect, as well presidential and vice presidential candidates. Twenty-First Amendment. The Twenty-First is the only Amendment that repeals, or cancels/invalidates, a previous Amendment. The Twenty-First repealed the Eighteenth Amendment and ended Prohibition at a national level, although it exclusively gave the States the power to prohibit the sale, manufacture or distribution of intoxicating liquors within their respective borders. Twenty-Second Amendment. George Washington set a standard for all future Presidents when he declined to seek a third term. This standard was either voluntarily followed by the President, or enforced by the voters, in every case until Franklin Roosevelt, who was elected four times. In order to restore Washington's tradition, the Amendment limited all future Presidents to a maximum of two terms. Twenty-Third Amendment. The Twenty-third Amendment gave citizens of the District of Columbia the right to vote in Presidential election. Twenty-Fourth Amendment. The Twenty-fourth prohibited the denial of a vote based on failure to pay a tax such as the poll tax. Southern states had used the poll tax to deny poor African-Americans the right to vote. Twenty-Fifth Amendment. The Twenty-fifth provided special provisions for an emergency such as a disabled but still living President. The Amendment allowed the Vice President and a majority of the Cabinet to formally declare the President unable to carry out his duties. If this was the case, then the Vice President would become Acting President until the President declared himself fit to continue. If the Vice President still felt that the President was not capable of continuing, then he and the Cabinet could again declare the President's disability. If this was the case, then Congress had to assemble to decide the matter. The President would continue unless the Congress, by a vote of two-thirds of each house, agreed with the Vice President and Cabinet. Twenty-Sixth Amendment. The Twenty-sixth provided that all persons eighteen years or older could not be denied the right to vote due to age. Twenty-Seventh Amendment. The Twenty-seventh Amendment was originally proposed by Congress in 1789 at the same time as the Bill of Rights. However, it was not ratified by legislatures of the required three-fourths of the states until 1992. It provides that Congressional salary changes cannot take effect until an election (of Representatives) occurs. 

Congress. As has been stated earlier, Congress includes two distinct houses- the Senate and the House of Representatives. The House of Representatives represents the people of the States based upon the population of each, while the Senate allows two Senators to each state regardless of population. The Legislative. The Congress is the Legislative Branch. Its main function is to make laws. It also oversees the execution of these laws, and checks various executive and judicial powers. The Congress is "bicameral"- it is composed of two houses. One house is the House of Representatives, the other is the Senate. The House of Representatives, or House for short, is currently composed of four hundred and thirty-five members. Each of the fifty states is allocated one or more representatives based on its population as calculated by the decennial (once in ten years) census. Each state is guaranteed at least one representative. A state that is allocated more than one representative divides itself, as state procedures dictate, into a number of districts equal to the number of representatives to which it is entitled. The people of each district vote to elect one representative to Congress (States that have only one representative allocated choose "at-large" representatives- the state votes as one entire district). The District of Columbia and a number of U.S. territories have been permitted to elect delegates to the House. These delegates may participate in debates, and sit and vote in committee, but are not allowed to vote in the full House. Every House member faces re-election in an even-numbered year and is elected to a two-year term. The House is presided over by a Speaker, who is directly elected by the members of the House. The Senate is the upper house of the United States' legislative branch, possessing only one hundred members to the house's four hundred thirty-five. Each state chooses two senators, regardless of that state's population. The Constitution originally dictated that a state's senators were to be chosen by the state's legislature; after the Seventeenth Amendment was ratified in 1913, senators were elected directly by the state's population. In contrast to the House's two-year terms, Senators are elected to a six-year stint in office. In addition, only one-third of the Senate stands for election during an even year. These differences between the two houses were deliberately put into place by the Founding Fathers; the Senate was intended to be a more stable, austere body, whereas the House would be more responsive to the people's will. The Vice-President is President of the Senate, but he/she only votes if there is a tie. The Senate also chooses a President Pro Tempore to preside in the Vice-President's absence (though, in practice, most of the time, senators from the majority take turns presiding for short periods). The Senate and the House are both required to approve legislation before it becomes a law. The two houses are equal in legislative power, but revenue bills (bills relating to taxation) may only originate in the House. However, as with any other bill, the Senate's approval is still required, and the Senate may amend such bills. The Senate holds additional powers relating to treaties and the appointments of executive and judicial officials. This power is known as "advice and consent." The Senate's advice and consent is required for the President to appoint judges and many executive officers, and also to ratify treaties. To grant advice and consent on treaties, two-thirds of the Senators must concur (agree). While most votes require a simple majority to pass, it sometimes takes three-fifths of senators to bring a bill to a vote. This is because Senate rules hold that a bill cannot be voted on as long as it is being debated--and there is no limit on how long a senator may debate a bill. Senators sometimes use this rule to filibuster a bill--that is, continue debating a bill endlessly so that it cannot be voted on. The only way to end a filibuster is for three-fifths of all Senators to vote for a cloture resolution, which ends all debate and brings the bill up for voting. Use of the filibuster tends to be controversial. Whichever party is in the majority tends to call its use "obstructionism," while the other side sees it as an important check on the majority. The House has the sole power to impeach federal executive and judicial officers. According to the Constitution, officers may be impeached for "treason, bribery, or other high crimes and misdemeanors.” The Senate has the sole power to try all such impeachments, a two-thirds vote being required for conviction. The Constitution requires that any individual convicted by the Senate to be removed from office. The Senate also has the power to bar that individual from further federal office. The Senate may not impose any further punishment, although the parties are still subject to trial in the courts. As the Vice-President (being next-in-line to the Presidency) would have an obvious conflict of interest in presiding at a trial of the President, in such cases, the Chief Justice presides. Interestingly, no similar provision prevents the Vice-President from presiding at his or her own trial. House of Representatives. Every ten years, the United States conducts a Census to determine the population of each state. Then, each state is apportioned, or allocated, a number of seats based on its Census population, with the more populous states receiving more seats than the less populous ones. Currently, one seat equals roughly 600,000 constituents. However, no matter how low a state's population is, it is always entitled to at least one seat. The state then conducts a process known as redistricting. In this process, the state divides itself into a number of districts of equal population; each district may then elect one representative. Of course, in states entitled to only one representative, the entire state is one district and redistricting does not occur. (See gerrymandering) In addition to representatives for the states, five American territories- the District of Columbia, Puerto Rico, Guam, the US Virgin Islands, and American Samoa- each choose non-voting members. Puerto Rico chooses a "Resident Commissioner" for a four-year term, while the others choose "Delegates" for two-year terms. Representatives hold office for two-year terms. The election for Representatives is held on the Tuesday immediately after the first Monday in November of every even numbered year. A Representative actually takes office on January 3rd after the election. No person may be a Representative unless he qualifies under the Constitution as follows. The requirements are: at least twenty-five years of age, inhabitance in the state of election, and citizenship of the United States for at least seven years. Though no specific number of Representatives is set by the Constitution, the law of the United States sets the number of Representatives at 435. The House of Representatives elects a Speaker who presides over the House. The Speaker is traditionally the leader of the majority party of the House. The Speaker of the House does have a significant role outside of the Congress in that she or he is third in line to the Presidency. Because of its size, the House relies heavily upon fixed rules and strict timetables for debate. When bills are debated on the floor of the House, each party's leader is allocated a fixed amount of time to present their argument for or against the bill, and they can appropriate this time to members of their party as they see fit. During House debates, it is common for representatives to "yield their time" to one another. Times for debate and other procedures are set by the House Rules Committee, which is generally considered to be one of the most powerful committees in Congress. Senate. Each state is entitled to two senators regardless of population. Originally, the state legislatures chose the senators. However, the people of the states choose their own senators at present. Senators hold office for six-year terms. Elections are held every two years, at the same time as the election for Representatives. The Senators are classified into three separate classes. At each election, the Senate seats of one particular class are up for election. Constitutional requirements for Senators are slightly more strict than those for Representatives. The qualifications are: at least thirty years of age, inhabitance in the state of election, and citizenship of the United States for at least nine years. The Vice President of the United States presides over the Senate and holds the title of President of the Senate. His power is considerably lesser than that of the Speaker of the House. He also does not have a vote in the Senate, unless the Senate is tied. As the Vice President normally does not attend unless there is a likelihood of a tie vote, the Senate chooses one of its members to be President of the Senate "pro tempore", or temporarily. The President pro tem, as he is often called, is normally the most senior Senator of the majority party. Just as in the house, the President or President pro tem does not preside during most meetings; this task is often given to new Senators so that they may learn the procedures of the body. This is not as easily possible in the House because of the much greater authority of whoever presides over the Representatives. In comparison to the House, the Senate has relatively few procedural rules, and no fixed schedules for debate. It is possible for a Senator to continue speaking for hours on end to delay unwanted legislation: this tactic is called a "filibuster." Any individual Senator is also allowed, by Senate rules, to stop the introduction of a bill with a motion from the floor, although this is almost never done in practice. Party leaders, committees, and caucuses. Both major political parties (the Republican Party and the Democratic Party) have designated floor leaders in both houses of Congress. The floor leader for the majority party is called the House (or Senate) Majority Leader, while the floor leader for the minority party is called the House (or Senate) Minority Leader. The second-in-command of each party's delegation is called the Whip, as their job is to "whip" other members of the party into action on various legislative measures. Each party's leadership is responsible for allocating its members to committees. There are a number of "standing committees" in each house, dedicated to various government functions such as the armed forces, education, and transportation. At any time, there are also several "select committees" that are set up for more timely problems such as government reforms. Occasionally, both houses of Congress will establish a "joint committee" to deal with certain issues. There are also many less formal associations in Congress, known as "caucuses," which are formed by members interested in various issues, such as relations with specific countries, ethnic issues, and industrial sectors. 

The Executive Branch, which executes and enforces the laws, is headed by the President and the Vice President. In addition, it includes the executive departments, which deal with general topics, and the heads of departments, who are known as Secretaries (Attorney-General in the Department of Justice). Each Department, in turn, is divided into a number of bodies, which are known as agencies, services, commissions, councils, bureaus, authorities, offices, administrations, and boards. The President. The President is the elected head of state and head of government of the United States. The President leads the Executive Branch and is the commander-in-chief of the United States Armed Forces. Article Two of US Constitution vest the executive power in the President. The power includes execution of federal law, alongside the responsibility of appointing federal executive, diplomatic, regulatory and judicial officers and making treaties with foreign nations, however all treaties made by the President must be approved by a two-thirds vote of the Senate. The president is further empowered to grant federal pardons and reprieves, and to convene and adjourn either or both houses of Congress under extraordinary circumstances. The President is also largely responsible for dictating the legislative agenda of the party to which he is enrolled. The President of the United States is often considered one of the most powerful people in the world. Powers and Duties. Article One: Legislative Powers. The first power the President is given in the Constitution, is the power to veto bills of legislation. The Constitution requires that all bills passed by Congress to be presented to the President before it can obtain legal force. Once a bill has been presented, the President has one of three options: Article Two: Executive Powers. One of the most important Presidential powers is the command-in-chief of the Armed Forces of the United States. While the Constitution vest the power to declare war solely in Congress, the President is ultimately reasonable for the disposition and direction of the Armed Forces. Along with the armed forces, the president also directs U.S. foreign policy. Through the Department of State and the Department of Defense, the president is responsible for the protection of Americans abroad and of foreign nationals in the United States. The president decides whether to recognize new nations and new governments, and negotiates treaties with other nations, which become binding on the United States when approved by two-thirds vote of the Senate.["citation needed"] Although not constitutionally provided, presidents also sometimes employ "executive agreements" in foreign relations. These agreements frequently regard administrative policy choices germane to executive power; for example, the extent to which either country presents an armed presence in a given area, how each country will enforce copyright treaties, or how each country will process foreign mail. However, the 20th century witnessed a vast expansion of the use of executive agreements, and critics have challenged the extent of that use as supplanting the treaty process and removing constitutionally prescribed checks and balances over the executive in foreign relations. Supporters counter that the agreements offer a pragmatic solution when the need for swift, secret, and/or concerted action arises The President is the head/leader of the Executive Branch of the federal government and is Constitutionally bound to "take care that the law be faithfully executed." The President has the power to make appointments within the Executive branch. Ambassadors, members of the Presidential Cabinet, and federal court officers are appointed by the President, although every appointed official by him must be approved by the Senate. The President also has the power to fire/remove purely executive officials at will, although Congress has the authority to curtail or constrain the Presidents power to fire/remove commissioners of independent regulatory agencies and certain inferior executive officers by statute. The president possesses the ability to direct much of the executive branch through executive orders that are grounded in federal law or constitutionally granted executive power. Executive orders are reviewable by federal courts and can be superseded by federal legislation. The president also has the power to nominate federal judges, including members of the United States courts of appeals and the Supreme Court of the United States. However, these nominations do require Senate confirmation. Securing Senate approval can provide a major obstacle for presidents who wish to orient the federal judiciary toward a particular ideological stance. When nominating judges to U.S. district courts, presidents often respect the long-standing tradition of Senatorial courtesy. Presidents may also grant pardons and reprieves for offense recognizable under federal law. The President cannot, however, grant pardons in impeachment cases, or grant pardons for persons convicted under State law/Jurisdiction. The President also cannot be simultaneously a member of Congress, neither can any other executive officer. Therefore, the President cannot directly introduce legislative proposals for consideration in Congress. However, the President can take an indirect role in shaping legislation, especially if the President's political party has a majority in one or both houses of Congress. For example, the President or other officials of the executive branch may draft legislation and then ask senators or representatives to introduce these drafts into Congress. The President can further influence the legislative branch through constitutionally mandated, periodic reports to Congress. These reports may be either written or oral, but today are given as the State of the Union address, which often outlines the President's legislative proposals for the coming year. Additionally, the President may attempt to have Congress alter proposed legislation by threatening to veto that legislation unless requested changes are made. Succession. If the President dies, resigns, is removed from office, or is unable to discharge the powers and duties of the Presidency, the Vice-President becomes President. Under the Twenty-Fifth amendment, the President may declare himself disabled and upon doing so, transfers the powers and duties of the Presidency to the Vice-President, who then becomes acting President. The President may resume the powers and duties of the Presidency by declaring himself again able to discharge the same. Under the same amendment, the Vice-President, along with a majority of the members of the President's cabinet, may declare the President disabled by submitting a written declaration to the Speaker of the House and the president "pro tempore" of the Senate, stating that President is disabled from discharging the powers and duties of the Presidency, and upon doing so, transfer the Presidential powers to the Vice-President as acting President. The President may resume discharging the powers and duties of the Presidency by submitting a written declaration to the Speaker of the House of Representatives and President "pro tempore" of the senate stating that such disability does not exist. If the Vice-President and Cabinet contest this claim, it is up to Congress, which must meet within two days if not already in session, to decide the merit of the claim. A President is: A person may not hold the office of President, if that person: The Vice President. The Vice President's only executive function in the Constitution is to become President in the event that the President dies or is incapacitated. The 25th Amendment provides that this may occur when: Additionally, "whenever there is a vacancy in the office of the Vice President, the President shall nominate a Vice President who shall take office upon confirmation by a majority vote of both Houses of Congress." In practice, the Vice President often takes an active role in policy making. The Vice President also serves as president of the Senate, but cannot vote except to break a tie. Elections. Each state is entitled to choose a number of Electors equal to the number of Representatives and Senators it elects to Congress. Thus, a state will choose no less than three electors- two for the Senators, and one for the minimum of one Representative. The state may choose its Electors in any way it pleases. However, all the states allow the people to choose all Electors. Forty-eight states use the following system: each candidate nominates a panel of electors. When a voter votes for a candidate, they actually vote for the nominated panel. Then, the candidate who wins more votes than any other candidate has his nominated panel appointed. Thus, a candidate who does not necessarily win all the votes will receive all of the state's electors. The exceptions to this system are Maine and Nebraska. These states allow a candidate to nominate one state panel of two electors, plus one elector for each district. Then, the candidate who wins the state has his two-member state panel appointed as electors, while the candidate who wins an individual district has his nominee for that district appointed. In all cases, the election for the Electors is held on the same day as Congressional elections. Regardless of how they are chosen, the Electors all assemble in their own states in December. The manner of the voting by Electors has changed. Firstly, the following is the original system: If, however, no candidate received a majority, then the House of Representatives would choose one of the top five candidates as President. In voting for President, each state casts one block vote. In any case, whichever other candidate holds the greatest number of votes other than the candidate elected President becomes Vice President. In case of a tie for second-place, the Senate elects the Vice President. The District of Columbia chooses Electors like the other states, but the District can in no event choose more Electors than any other state. Executive departments. At present, there are fifteen executive departments. They are the departments of: Each department is subdivided into a number of agencies, bureaus and other divisions. Executive agencies. In addition to the Cabinet departments, there are certain independent bodies which are not part of any department, but report directly to the Executive Office of the President. These include: The Cabinet. The Cabinet currently consists of the Vice President, the White House Chief of Staff, the heads of each executive department, and the heads of the EPA, OMB, ONDCP and USTR. 

The Courts. The United States judicial system includes the Supreme Court of the United States and the inferior federal courts. The President nominates an individual to serve as a judge, after which the Senate must grant its advice and consent before the President can formally appoint the judge. A judge holds office during "good behavior", which is usually interpreted as meaning a life term. The Supreme Court. The Constitution creates the Supreme Court, but permits Congress to set the number of Justices. Currently, the Supreme Court includes the Chief Justice of the United States and eight Associate Justices. The Supreme Court has original jurisdiction over limited categories of cases, such as cases between two or more states. It hears most of its cases through its appellate jurisdiction, in which it hears appeals from lower federal and state courts. It exercises its appellate jurisdiction selectively; four of the nine Justices must grant a "writ of certiorari" before a case can be heard. The Court of Appeals. The United States is divided into twelve regional "circuits", each of which has a Court of Appeal. Additionally, there is a Federal Circuit which hears appeals from certain special tribunals and courts. The regional circuits (which are officially known by a number only, except for the DC circuit) are as follows: Each Court of Appeal includes a different number of members. The First Circuit has the fewest with six members, while the Ninth Circuit has the most with twenty-eight. District Court. Every state is divided into one or more Court Districts (which are distinct from Congressional districts). The total number of districts is ninety-four. Each district court includes a different number of members. The Eastern and Western District of Kentucky, the Eastern District of Oklahoma, and the Northern, Eastern, and Western District of Oklahoma, each have the fewest judges with one, while the Central District of California has the most with twenty-seven. Other Courts. The Congress has established special bankruptcy courts and other courts to rule on specific matters. 

The Japanese language uses post-position particles (助詞; じょし) to denote the direction of an action and who is performing the action. They consistently come after the word that they modify. There are three particles used very frequently in the language: は, を and が. This module covers these along with a few other common ones but an exhaustive list would run very long. The topic and subject markers は and が. The particle "は" (pronounced as "わ" when used as a particle) is the "topic marker" denoting topic of discussion, while "が" is the "subject marker" and marks a noun that performs an action. The difference between the two tends to cause confusion among beginners but their usage can be summed up as matter of focus. The topic particle "は" is used when introducing a topic and gives focus to the "action" of the sentence (i.e., the verb or the adjective). The subject marker "が" is used when emphasising the subject giving focus to the "subject" of the action. One can also think of it as replacing "~は" with the phrase "as for ~", "on the topic of ~" or "regarding ~" to distinguish it from "が". While these phrases aren't common in English we can use these expressions here to better show the main difference between "は" and "が". The difference can also be displayed by using both subject and topic markers in one sentence: One has to be careful using both "は" and "が" in one sentence. If a verb is actually acting on the (direct) subject, usually a different particle (like を) has to be used. "は" is generally more flexible, because the "it" can be assumed, and is therefore recommended to novices who have not grasped the difference between the two. "は" also has the specialized function of being used for comparisons as well. Often the grammatical subject may also be the topic. In this case, "は" normally replaces "が". However, if the subject is never known, you cannot use "は" and must use "が". This is similar to using pronouns: You can't state, "It is over there", without first stating what "it" may be. The direct object marker を. The particle "を" (predominantly pronounced "お") is the "direct object marker" and marks the recipient of an action. It also indicates the place through which the action occurs: As with much of the language, parts of a sentence that can be assumed from context are often omitted and the direct object particle is commonly dropped in conversational (colloquial) Japanese. O is commonly used to identify the object in which the verb is affecting. for example we will use the sentence ( I drink juice ) ( Watashi wa juice o nomu ) o is identifying the word Juice as the object in which nomu's action is taking place. Nomu in Japanese means Drink / to drink. O is basically telling us that the word juice is the object that the verb is interacting with. The indirect object marker に. "に" marks the verb's "indirect object" (i.e. the destination of a targeted verb action) translating as "to", "in", "at" or "by". It also indicates the location touched or affected by an event or action: "に" can also be used as an "object of a preposition" marker when found in prepositional phrases like の前に (no mae ni), which means "in front of" or "before" depending on the context of the sentence. The particle "へ" described below is used exclusively for marking the destination. The destination marker へ. へ (pronounced "え" when used as a particle) indicates the direction of an action, roughly the equivalent of "to" or "toward" in English. The question marker か. Placing か at the end of a sentence changes a statement into a question. Use it at the end of a verb to make it a question, or at the end of an interrogative pro-form to make it into a demonstrative pronoun. For more on the question marker, see: ../Sentence ending particles/. The possessive marker の. "の", is most commonly used as a "possessive marker" (similar to the English " 's "). The particle can also function as a noun link, indicating that the preceding noun (or adjectival noun) modifies the following noun. It can also be used for "nominalisation", converting verbs and (proper) adjectives into nouns. Note that in this last example two particles are used together: の and が: the first makes the action a noun, and the second tells that this action is what the sentence is all about. The exhaustive list conjunction と. This particle acts as a conjunction on the words it separates. Unlike conjunctions of more than two words in English, where only the last two are separated with an "and" and the rest with commas, the Japanese conjunction separates each word and commas are not used. This applies to exhaustive lists, i.e. when all objects are explicitly mentioned. The particle is used to indicate parallelism with the subject, often meaning "with": The incomplete list marker や. This particle is used to connect various words implying that the listing is not exhaustive. The particle "など" may be added after the list to emphasise that the list is incomplete. The "also" marker も. も is quite simply a marker that says "also". It replaces the particles は, が and を but can also follow other particles. This can also be used to form a large list of words all acting as though one of the basic particles (は, を, or が) were affecting the whole list. Worth noting is that used with an interrogative pro-form (e.g. who, where, how) the も particle negates the pro-form: The means particle で. The particle で can be used in several situations indicating means. These can be for example an instrument, a location or a language. As a note of interest, the で from the copula である is also actually an instrumental-maker. で marks the whole previous expression instrumental to the verb ある. However, this is the classical meaning of the copula and rarely "explicitly" treated this way in modern Japanese. Origin and limit から and まで. These particles indicate the starting point or border of an action. This may be a location as well as a time and corresponds roughly with "from" and "until". 

"The Once and Future King" is a novel by T. H. White about the legend of King Arthur. It is often assigned reading in English literature classes and is composed of five books: While the first four were originally published separately, and reworked for inclusion in "The Once and Future King", the last was to be published in the greater work for the first time. This was vetoed by White's publisher, supposedly due to wartime paper restrictions. It wasn't until after White's death that The Book of Merlyn joined the other four. 

"The Once and Future King" - The Sword in the Stone Book I covers the Wart's (Arthur's) childhood education by Merlyn, including his adventures into Morgan le Fay's castle with Robin Wood, and his transformations into many different animals. It ends with the Wart removing the sword from the stone after King Uther's death, and his subsequent coronation 

The Once and Future King - Book I: The Sword in the Stone - Chapter 17 Plot summary. The chapter begins a few months after the boar-hunt. King Pellinore gives the Questing Beast, now healthy, a two-hour headstart in the chase. Merlyn, Archimedes, and the Wart discuss the birds, a subject that has arisen due to spring's arrival. The Wart asks to become a bird again, because the last time he didn't get a chance to fly; Archimedes plans to go with him at night. The Wart then mentions his favorite bird is the rook, because rooks fly just for pleasure, not necessarily with a purpose. Archimedes mentions rooks are one of few birds with a legal and social system, and in response to Merlyn says his favorite are pigeons. Then Merlyn says he prefers the chaffinch, and continues that bird calls arise out of imitation of their prey or their environment. Finally Kay enters, ironically saying that he was shooting birds. Analysis. Anachronisms - Merlyn mentions both and , and Archimedes compares pigeons to . All came much later. Animals as forms of government - Archimedes mentions that the rooks have "parliaments...and a social system;" (a play upon the for rooks being a "parliament") they make "laws about the defense of the rookery, and marriage, and so forth" 

Under the Kingdom of Great Britain, the American colonies experienced five situations which would guide them in creating a constitution. The British Parliament believed that it had the right to impose taxes on the colonists; it had "virtual" representation over the entire empire, while the colonists believed Parliament had no such right, as they had no "direct" representation in Parliament. By the 1720s all but two of the colonies had a locally elected legislature and a British appointed governor. Often, these two branches of government would clash, with the legislatures imposing their "power of the purse" to control the British governor. Thus, Americans viewed their legislative branch as a guardian of their liberty, while the executive branches was deemed tyrannical. There were several examples of royal actions that upset the Americans. For example, taxes on the importation of lead, paint, tea, paper, spirits, rum, wine, molasses, sugar, and other products were imposed at various times. Also, the Parliament provided for a duty to be paid on court documents, certificates, licenses, deeds, other legal documents, playing cards, pamphlets, books, calendars, newspapers, and other papers, as well as dice. The variety of taxes imposed, as well as other causes, led to the Americans' disdain for the British system of government. After the Boston Tea Party, the Parliament of Great Britain and the King passed Acts that outlawed the Massachusetts legislature. The Parliament also provided for special courts in which British judges, rather than American juries, would try colonists. The Quartering Act and the Intolerable Acts required Americans to provide room and board for British soldiers. Americans especially feared British actions in Canada, where civil law was once suspended in favor of British military rule. American distaste for the system of British government would lead to revolution. Americans had formed their own local institutions which were not British at all, but American. The political ideas of the Americans actually had their root in the British radicals of the early 18th century. England had passed beyond those ideas by 1776 and the resulting conflict resulted in the first American attempts at a national government. 

The Articles of Confederation. The Articles of Confederation accomplished certain things, but without a strict leader, or a government that could really do anything to help they turned out to be a bad thing for the United States at that time. First, they expressly provided that the states were sovereign. (A "sovereign state" is a state that is both self-governing and independent.) The United States as a Confederation was much like the present-day European Union. Each member was able to make its own laws; the entire Union was merely for the purposes of common defense. The reason for the independence of the colonies is clear- the colonies were afraid of the power of a central government such as the one in the State of Great Britain. The Articles provided that a Congress, consisting of two to seven members per state, would hold legislative power. The states, regardless of the number of Congress members representing them, each have one total vote. The Congress was empowered to settle boundary and other disputes between states. It could also establish courts with jurisdiction over the seas. Also, it could tax the states, even though it did not possess the power to require the collection these taxes by law. Faults of the Articles. The Congress, overall, was absolutely ineffectual. The Congress had to rely on the states for its funding. Since it could not forcibly collect taxes, the states could grant or withhold money and force Congress to accept their demands. Because it could not collect taxes, Congress printed paper dollars. This policy, however, absolutely wrecked the economy because of an overabundance of paper dollars, which had lost almost all value. The several states also printed their own currency. This led to much confusion relating to exchange rates and trade; some states accepted the currency of others, while other states refused to honor bills issued by their counterparts. Furthermore, the Articles included certain fallacies. For instance, it suggested that the approval of "nine states" was required to make certain laws. However, it made no provision for additional states. Thus, it would appear that the number nine would be in effect even if that number would actually be a minority of states. Also, the Articles required the approval of all states for certain important decisions such as making Amendments. As the number of 'States" would grow, securing this approval would become more and more difficult. The Conference at Annapolis. At Annapolis, Maryland, the delegates of the thirteen states were supposed to meet to discuss various changes to the Articles to grant more authority to Congress. However, eight of the states failed to send representatives. Thus, the Conference did not even occur. However, another Conference was called for 1787. This Conference at Philadelphia is what we now know to be the Constitutional Convention. 

The Legislature. The United States were basically divided into two classes- the large (more populous) states and the small (less populous) states. The large states included Pennsylvania, Virginia, and Massachusetts. The small states included Rhode Island, New Jersey, Delaware, Connecticut, New Hampshire, and even Maryland. Also, one may consider Georgia and the two Carolinas as small states, but these states hoped to increase their population and become large by importing slaves and attracting "immigrants" from other states. These were called “in-between states” The large states wanted to have proportional representation in Congress. They wished that the more populous states have more representatives than the less populous states. However, fearing that they would be overwhelmed by large numbers of representatives from other states, the small state delegates suggested that all states receive equal representation like under the Articles. James Madison of Virginia proposed a plan, which was presented by Edmund Randolph, supported by the large states, the Virginia Plan. It entailed: Meanwhile, New Jersey politician William Paterson proposed a plan on behalf of the small states. It involved: Thirdly, Alexander Hamilton of New York proposed a plan extremely similar to the British government. The British plan included: Hamilton's plan was rejected very quickly- it reminded the delegates too much of the tyranny and unhappiness under the King of the State of Great Britain. Connecticut Delegate Roger Sherman suggested that the small and large states compromise. He felt that the large states would never accept equal representation, while the small ones would never accept just proportional representation. His compromise, known as the Great Compromise, suggested the following: Though Sherman's compromise was initially rejected, the delegates were forced to accept it eventually. Otherwise, the Convention would have clearly broken down on the issue of representation. The Executive. Once the issue of representation was resolved, other issues seemed relatively easy to negotiate. The delegates continued to compromise on several issues, including the executive. Firstly, the delegates were concerned about a single individual as executive. The King, they said, was an individual with too much power. However, the argument failed when some pointed out that every single state in the union had one Chief Executive called a President or a Governor, rather than a Council of Presidents or Governors, and none of the states suffered from that Governor's tyranny. Similarly, the executive was granted substantial but not absolute power, after the example of the individual states. The manner of choosing the executive was the only one of concern. The following were proposed as electors for the President: The Framers rejected the idea of election by the People because they felt that, it would be impractical in the days of difficult communication, and inappropriate because the people would "naturally" vote for local candidates without any regard for those from other states. Also, they rejected the state or Congressional choice because they assumed that the President would feel indebted to and controlled by the states or the Congress. Such a problem would be present with any permanent body. Thus, they established a temporary body whose sole purpose was to elect the President- the Electoral College. (See Part III, Chapter 2.) Slavery. The problem of slavery, after the issue of representation, was probably the most dangerous one for the Convention to tackle. If the Convention adopted a plan that upset one region, then the states of that region might have withdrawn from the Convention, breaking up the meeting. Related to the issue of representation was the counting of slaves to decide the population of a state for the purpose of proportional representation in Congress. The South wanted slaves to count, but the North feared that the South could increase its power in Congress by importing more slaves. The Three-Fifths Compromise suggested the same standard as the Article of Confederation-"other persons," or slaves would be counted as three-fifths of persons. The three-fifths rule would be applied for deciding proportions in Congress and amounts of direct tax due from each state. Another compromise relating to slavery involved the importation of slaves. The Constitutional Convention compromised by allowing the slave trade to continue until 1808, when the Congress could lawfully ban it. Conclusion. The tired delegates were faced with a problem - that of a Bill of Rights. The delegates, however, refused to take the risk of breaking up the Convention and wasting hard work by debating specific rights. Thus, they assumed that a newly assembled Congress would add these Amendments, or they felt that the present Constitutional protections were sufficient. In order for the Constitution to gain effect, the Convention required that nine states approve it. In addition, the states not ratifying, or approving, the Constitution would not be subject to it. 

Introduction. The Constitution required ratification by nine states in order to come into effect. The fight for ratification was long and difficult. An important factor in ratifying the constitution was that the delegates to the Constitutional Convention required that ratification be done by special ratifying conventions, not by state legislature. Being interested in retaining state powers, the states would understandably have been resistant to ratifying a new, stronger central government. Federalists versus Anti- Federalists. Those who favored ratification were known as Federalists,while those who opposed it were considered Anti- Federalists. The Federalists attacked the weaknesses of the Articles of Confederation. They acknowledged that the Constitution was not perfect, but they said that it was much better than any other proposal then made. Three Federalists- Alexander Hamilton, James Madison, and John Jay- wrote a series of essays called "The Federalist Papers". The essays explained the constitution and defended its provisions. The documents were intended for the state of New York, though people from across the country read them. The Federalists defended the weakest point of the Constitution- a lack of a Bill of Rights- by suggesting that current protections were sufficient and that the Congress could always propose Amendments. Anti- Federalists such as Patrick Henry attacked the Constitution, suggesting that it would lead to a dangerously powerful national government. The Anti- Federalist arguments relating to the Bill of Rights were especially powerful. However, many were convinced by the assurances provided by the Federalists. State Conventions. Each state was to hold a convention to debate the Constitution and ratify or reject it. The Constitution was proposed in September, 1787. By the end of the year, some states that were in favor of the document ratified. Delaware's convention approved it by a vote of 30- 0; Pennsylvania's by a vote of 46- 23; New Jersey's by a vote of 38- 0. The next year, Georgia ratified by a vote of 26- 0 and Connecticut followed with a vote of 128- 40. The Federalists were already more than halfway to the nine-state margin. However, the states that did not ratify the Constitution in that year included the extremely important states of Massachusetts, New York, and Virginia. Massachusetts ratified the document by a close margin (187- 168) in February, 1788. Maryland followed with a 63- 11 vote, and South Carolina did the same with a 149- 73 vote. Then, New Hampshire provided the all-important ninth ratification by a 57- 47 vote. The United States was established under the new Constitution, but the important commercial state of New York and the populous state of Virginia, among others, still did not ratify. After tough battles, these states also ratified, Virginia by a 89- 79 margin and New York with 30- 27. The Bill of Rights was then created under the new Constitution, leading to North Carolina's ratification by a vote of 194- 77. Seeing itself alone, Rhode Island finally agreed to ratify with a 34- 32 vote. All thirteen colonies had ratified the Constitution by May, 1790. 

High SchoolPure Maths Extensions Introduction. This online textbook is intended for, but not limited to, high school students without a rigorous understanding and knowledge of university-level mathematics. Therefore, the text's language reflects the expected mathematical maturity of the intended audience. This book introduces several interesting topics not covered in the standard high school curriculum of most countries. The materials presented can be challenging, but at the same time, we strive to make this book readable to all who are a few years from applying to higher education.  From the authors It is our firm belief that math textbooks should not just be a collection of mathematical facts carefully laid out for rote memorization and cram sessions. A math textbook, especially for the youth, should be full of questions, not just exercises. These questions require some thought to answer and spark curiosity. After all, the questions keep the students engaged, not the answers. We sincerely hope to interest, stimulate, and challenge all those who read this book. Authors &amp; Contributors. A number of persons not listed below have also made important contributions to this book. Contributors are encouraged to edit and include themselves in this list. 

Slovio is a newly-constructed language created by linguist Mark Hučko. It is an international auxiliary language primarily created to help Slavic speakers communicate. The grammar of Slovio is similar to Esperanto, but the vocabulary is derived from the most common words from Slavic languages, that of the largest European-language group. The name of the language, Slovio, comes from the old-Slavic word "slovo" which means "word". 

Introduction. We have already considered moduli and modular arithmetic back in ../Number theory/, however in this section we will take a more in depth view of modular arithmetic. For revision, you should review the material in number theory if you choose. Simultaneous equations. When we speak of simultaneous equations with relation to modular arithmetic, we are talking about simultaneous solutions to sets of equations in the form There are two principal methods we will consider, "successive substitution" and the "Chinese remainder theorem". Successive substitution. The method of successive substitution is that where we use the definition of the modulus to rewrite these simultaneous equations, and then successively make substitutions. It will probably be best to motivate the idea with an example. Example: Solve 3"x" ≡ 10 (mod 19), and "x" ≡ 19 (mod 21) using successive substitution. First: Find the inverse of 3 in Z19; 3-1=-6, then Substitute in the second equation Find the inverse of 19 in Z21; 19-1=10 Writing in the equivalent form Substituting back j in (*) Writing back in the first form which is our solution. Chinese remainder theorem. The "Chinese remainder theorem" is a method for solving simultaneous linear congruences when the moduli are coprime. Given the equations multiply the moduli together, i.e. N=m1m2...mk, then write n1=N/m1, ..., nk=N/mk. We then set yi be the inverse of ni mod mi for all i, so yini=1 mod mi. Our solution will be To see why this works consider what values x mod mk takes. The term akyknk mod mk becomes equal to ak as yknk=1 mod mk, and all the terms ajyjnj mod mk become equal to zero as when formula_1 mk is a factor of nj. The Chinese Remainder Theorem is of immense practical use, as if we wish to solve an equation mod M for some large M, we can instead solve it mod p for every prime factor of M and use CRT to obtain a solution mod M. Powers and roots. This section deals with looking powers of numbers modulo some modulus. We look at efficient ways of calculating If we tried to calculate this normally - by calculating "a""b" and then taking the modulus - it would take an "exorbitant" amount of time. However some of the theory behind modular arithmetic allows us a few shortcuts. We will look at some of these and the theory involved with them. Fermat's (little) Theorem. Fermat's theorem allows us to see where "a""b" (mod "m") is 1. This has an application in disproving primality. It states So, for example, 1310=1 in Z11. Primitive elements. If in Zn, can we write some elements as powers of an element? This is conceivably possible. Let's look at Z3. The elements {1,2} constitute in fact :Z3*. Generally, we have Orders. We can express this idea in a different way, using the concept of the "order". We denote the order of "a" ∈ Zn* by the smallest integer "k" written On(a) such that For example, On(-1)=2 for all n except 2, since except when n = 2, since in that field -1 = 1 and thus has order 1. Note if gcd("a","n")≠1, that is, "a" ∉ Zn*, the order "is not defined". Properties of orders. The orders obey some properties, the first of which was originally proven by Lagrange: If p prime, gcd(a,p)=1, Orders and finding primitive elements. Given these facts above, we can find primitive elements in Zp for "p" &gt; 2 fairly easily. Using the above facts, we only need to check "a"("p"-1)/"p"i="x"i in Zp for all "i", where the "p"i are the prime factors of "p"-1. If any of the "x"i are 1, "a" is not a primitive element, if none are, it is. Example: Find a primitive element of Z11. Try 2. "p"-1 = 10 = 2 . 5 Check: Neither is 1, so we can say that 2 is a primitive element in Z11. Problem set. Given the above, answer the following. (Answers follow to even-numbered questions) Euler's totient function. Euler's totient function is a special function that allows us to generalize Fermat's little theorem above. It is defined as Some results. We have the following results leading on from previous definitions. In other symbols: formula_2. "Proof of 2.": There are "p"k elements in Z"p"k. The non-invertible elements in Z"p"k are the multiples of "p" and there are "p"k-1 of them: "p", 2"p", 3"p", ..., ("p"k-1-1)"p", "p"k. Removing the non-invertible elements from the invertible ones leaves "p"k-"p"k-1 left. ∎ "Corollary to 1, 2 and 3": If "n" has distinct prime factors (i.e. not counting powers) "p"i for i=1...,r we have For example: "Proof of 3.": We can prove this equality using a special case of the Chinese Remainder Theorem, where the CRT is now just a system of 2 congruences, namely: (remember that the CRT is applicable here because m and n are assumed coprime in the equality). Note that a1 can take on m values (from 0 to m-1), and a2 can take on n values (from 0 to n-1). Also note that, for each and everyone of the m*n (a1, a2) tuples, there is a unique solution x that is strictly smaller than m*n. Moreover, for each x strictly smaller than m*n, there is a unique tuple (a1, a2) verifying the congruence system (these two assertions are a component of the Chinese Remainder Theorem: a solution to the congruence system is unique modulo m*n). With this bijective uniqueness property in mind, the proof is simple. Go through each x, from 0 to m*n-1, and show that if x is a totient of m*n (i.e., gcd (x,m*n) = 1), then a1 is a totient of m and a2 is a totient of n. Furthermore, you must also show that if a1 and a2 are totients of m and n respectively, then it follows that x must be a totient of m*n. If gcd (x,m*n) = 1, then according to Bezout's identity, there exist X and Y integers such that x*X + m*n*Y = 1. Furthermore, we have: Therefore, a1*X + m*(k + n*Y) = 1, &lt;br&gt; should this be a1*X + m*(k*X + n*Y) = 1 ?? &lt;br&gt; so gcd (a1,m) = 1, and therefore a1 is a totient of m. Proceed similarly to prove that a2 is a totient of n. Proving the other direction is very similar in that it requires some simple replacement algebra. So what have we shown? In the above we have shown that for every totient x of m*n, there is a unique tuple of totients of m on the one hand and n on the other hand. Furthermore, that for each tuple of totients of m on the one hand and n on the other hand, there is a unique totient of m*n. Therefore, phi(m*n) = phi(m)*phi(n). "Proof of 4.": Let Q(g) be the set of all integers between 1 and n inclusive, such that gcd(x,n) = g. Q(g) is nonempty if and only if g divides n. If g doesn't divide n, then good luck finding an x such that g is the greatest common DIVISOR of x and n. Secondly, if x belongs to Q(g) for a given g, then it can't belong to another Q(...), since, if n is fixed, then gcd(x,n) is unique, by definition of the GREATEST common divisor. Thirdly, for all x between 1 and n inclusive, there exists a g such that gcd (x,n) = g (in the "worst" case, it's 1). Put together, these three properties imply that the union of all the Q(g) sets (for each g a divisor of n), which are pairwise mutually exclusive, is the set {1,2,3...,n}. And therefore, the sum of the cardinalities of each Q(g) equals n. Now we show that |Q(g)| = φ(n/g). One direction: Let x be an arbitrary member of Q(g) for some g. Therefore, we have that gcd (x,n) = g =&gt; gcd (x/g, n/g) = 1 =&gt; x/g belongs to the set of numbers coprime to n/g (whose cardinality of course is φ(n/g)). For diff\ erent x's, the two values x1/g and x2/g are distinct. So for each x in Q(g), there is a correspondingly unique x/g in the set of numbers coprime to n/g. Other direction: Let x be an arbitrary member of the set of numbers coprime to n/g. This implies gcd (x,n/g) = 1 =&gt; gcd (xg,n) = g =&gt; xg belongs to Q(g). For different x's, the two values x1g and x2g are distinct. So for each x in the set of numbers coprime to n/g, there is a correspondingly unique xg in Q(g). Therefore, |Q(g)| = φ(n/g). Euler's theorem. We can now generalize Fermat's theorem to extend past just Zn. Euler's theorem says: Example: Find 3216 in Z14. We need to calculate firstly φ(14)=φ(7)φ(2)=(7-1)(2-1)=6. Then write the exponent as: 216 = 6 × 36 So: 3216=(36)36 But Euler's theorem tells us 36=1 in Z14 (i.e., mod 14) since 3φ(14)=1 in Z14 as above. So we have: 3216=136=1. Calculating large powers efficiently. When Euler's or Fermat's theorem fails us in the calculation of a high power, there is a way to decompose an exponent down so calculation is still easy. Let us work through an example as motivation. Example. 528 in Z4. First write 28 in base 2 = (11100)2 = 24+23+22 = 16 + 8 + 4 Now 528 = 516+8+4 = 516 58 54 Now rewrite these powers of 2 as repeated exponents: When you calculate each exponent, reduce mod 4 each time. Problem set. Given the above, calculate the following powers. (Answers follow to even-numbered questions) 

Look at the diagram above. What is the potential difference across the 4Ω resistor and each of the 2Ω resistors ? 

I'm sorry but this is the wrong answer. The 4Ω resistor gets the full 6V and the two 2Ω resistors have to share the 6V between them because they are in series. They therefore get 3V each. On to the next question 

I'm sorry but this is the wrong answer. The 4Ω resistor gets the full 6V and the two 2Ω resistors have to share the 6V between them because they are in series. They therefore get 3V each. On to the next question 

I'm sorry but this is the wrong answer. The 4Ω resistor gets the full 6V and the two 2Ω resistors have to share the 6V between them because they are in series. They therefore get 3V each. On to the next question 

Well done! this is the correct answer. On to the next question 

This is the same circuit as in question 3. But this time you are being asked to work out the current through each resistor. 

Act: A Bill that has been passed by Congress and become law. (Note: "Act" refers to a bill that has become law. "Act of Congress" refers to a bill passed by Congress but has not yet become law. "Act of the Senate" or "Act of the House" refers to a bill passed by one house only.) Amendment: A change to a Constitution or to a bill. Anti-Federalist: One who opposed the ratification of the Constitution. Bill: Proposed legislation. Bicameral: Having two branches of a legislature. Bill of Rights: The first ten amendments to the United States Constitution. Checks and Balances: A system whereby the different branches of government balance each other so that one branch does not gain too much power. Commerce: Trade or exchange of goods and money. Concurrent Powers: Powers shared by both the federal and the state governments. Delegate: A representative entitled to exercise powers on behalf of a government, such as a in a convention. Equal Representation: System under which all political entities such as states receive representation in the legislature equal to each other. Executive Branch: The branch of government in charge of enforcing and executing the laws. Federalism: System under which a national government as well as regional governments (the states) have certain powers of legislation. Federalist: One who supported the ratification of the Constitution. Impeachment: An accusation made by a legislature, or part of legislature, against an executive or judicial officer. The Impeachment is only the accusation and does not indicate guilt, which is determined at a trial in the other part of the legislature. Judicial Branch/ Judiciary: The branch of government in charge of interpreting the laws; the courts. Judicial Review: The power of a court to rule laws unconstitutional and therefore null and void. Legislative Branch/ Legislature: The branch of government in charge of making the laws and overseeing their enforcement. Probable Cause: Cause (for an action such as searching a home) that has a reasonable basis. Proportional Representation: System under which a political entity such as a state receives representation in the legislature in proportion with its population. Quarter: To provide room and board for. Usually used in the context of boarding soldiers. Ratification: 1. the approval of a constitution or an amendment by a state through a legislature, convention, or other method. 2. the executive act of approving a treaty. In the US, the President may ratify treaties, but only with the advice and consent of two-thirds of the Senate. Subpoena: A court order commanding a person or entity either to surrender documents to the court or to testify. Veto: The rejection of a bill by the executive. 

The Once and Future King - Book I: The Sword in the Stone - Chapter 21 Plot summary. Wart is sad, because Kay is being made a knight, and he will be only Kay's squire. Merlyn suggests that the Wart learn something, mentioning on the side the worthlessness of education to a laborer. As he warns that this is the last use of his magic, he turns Wart into a badger. Wart picks on a hedgehog and threatens to eat it, and the hedgehog asks for mercy by singing. Finally, when the hedgehog reveals he is one that was in Merlyn's hut, the Wart agrees to let the hedgehog go. Then, the Wart goes to see Badger, who tells a myth about how, when God was about to shape the embryos of the species, most animals requested specialization of their parts to help them. Man, though, said he would settle for the embryo shape, because God must have some plan in His mind. After God reveals man has solved the riddle, he gives them dominion over the other animals. Badger continues, questioning Man's ability at this leadership; he mentions he can only think of seven species that war against themselves. When the Wart protests, he asks, "Which did you like best, the ants or the wild geese?" Analysis. Animals as governments - Badger ends by asking Wart whether he preferred the species that made war or the species that kept peace. Though Wart sees war as possibly noble, Badger reminds him that peace can accomplish much more. Badger also questions if humanity really holds the mankind is taken as an animal species. 

For a chapter: The Once and Future King - Book N: ../The Book Title/ - Chapter X Plot summary. ../Arthur/ does something. Analysis. Topic - For a book: "The Once and Future King" - Book N: The Book Title Book N covers topics with some characters. 

Introduction. Although the real numbers can, in some sense, represent any natural quantity, they are in another sense incomplete. We can write certain types of equations with real number coefficients which we desire to solve, but which have no real number solutions. The simplest example of this is the equation: Your high school math teacher may have told you that there is no solution to the above equation. He/she may have even emphasised that there is no "real" solution. But we can, in fact, extend our system of numbers to include the "complex" numbers by declaring the solution to that equation to exist, and giving it a name: the "imaginary unit", formula_2. Let's "imagine" for this chapter that formula_3 exists. Hence "x" = "i" is a solution to the above question, and formula_4. A valid question that one may ask is "Why?". Why is it important that we be able to solve these quadratics with this seemingly artificial construction? It is interesting delve a little further into the reason why this imaginary number was introduced in the first place - it turns out that there was a valid reason why mathematicians realized that such a construct was useful, and could provide deeper insight. The answer to the question lies not in the solution of quadratics, but rather in the solution of the intersection of a cubic and a line. The mathematician Cardano managed to come up with an ingenious method of solving cubics - much like the quadratic formula, there is also a formula that gives us the roots of cubic equations, although it is far more complicated. Essentially, we can express the solution of a cubic formula_5 in the form formula_6 An unsightly expression, indeed! You should be able to convince yourself that the line formula_7 must always hit the cubic formula_8. But try solving some equation where formula_9, and you run into a problem - the problem is that we are forced to deal with the square root of a negative number. But, we know that in fact there is a solution for x; for example, formula_10 has the solution x = 4. It became apparent to the mathematician Bombelli that there was some piece of the puzzle that was missing - something that explained how this seemingly perverse operation of taking a square root of a negative number would somehow simplify to a simple answer like 4. This was in fact the motivation for considering imaginary numbers, and opened up a fascinating area of mathematics. The topic of Complex numbers is very much concerned with this number "i". Since this number doesn't exist in this real world, and only lives in our imagination, we call it the "imaginary unit". (Note that formula_2 is not typically chosen as a variable name for this reason.) The imaginary unit. As mentioned above Let's compute a few more powers of "i": As you may see, there is a pattern to be found in this. Complex numbers as solutions to quadratic equations. Consider the quadratic equation: The "x" we get as a solution is what we call a complex number. (To be nitpicky, the solution set of this equation actually has two complex numbers in it; either is a valid value for x.) It consists of "two" parts: a "real" part of 3 and an "imaginary" part of formula_18. Let's call the real part "a" and the imaginary part "b"; then the sum formula_19 is a complex number. Notice that by merely defining the square root of negative one, we have already given ourselves the ability to assign a value to a much more complicated, and previously unsolvable, quadratic equation. It turns out that 'any' polynomial equation of degree formula_20 has exactly formula_20 zeroes if we allow complex numbers; this is called the Fundamental Theorem of Algebra. We denote the "real" part by "Re". E.g.: and the "imaginary" part by "Im". E.g.: Let's check to see whether formula_24 really is solution to the equation: Exercises. Substitute "z" and "w" into the quadratic equation above using the values you have computed in Exercise 3 and 4. What do you observe? What conclusion can you draw from this? Arithmetic with complex numbers. Addition and multiplication. Adding and multiplying two complex number together turns out to be quite straightforward. Let's illustrate with a few examples. Let "x" = 3 - 2"i" and "y" = 7 + 11"i", and we do addition first and now multiplication Let's summarise the results here. But how do we calculate: Note that the square root is only above the 5 and not the "i". This is a little bit tricky, and we shall cover it in the next section. Exercises:. Compute: Division. One way to calculate: is to rationalise the denominator: Utilising a similar idea, to calculate we "real"ise the denominator. The denominator is the sum of two squares. We get: If somehow we can always find a complex number whose product with the denominator is a real number, then it's easy to do divisions. If and Then "zw" is a real number. This is true for any 'a' and 'b' (provided they are real numbers). Exercises. Convince yourself that the product of "zw" is always a real number. Complex Conjugate. The exercise above leads to the idea of a complex conjugate. The complex conjugate of "a" + "ib" is "a" - "ib". For example, the conjugate of "2 + 3i" is "2 - 3i". It is a simple fact that the product of a complex number and its conjugate is always a real number. If "z" is a complex number then its conjugate is denoted by formula_41. Symbolically if then, The conjugate of "3 - 9i" is "3 + 9i". The conjugate of "100" is "100". The conjugate of "9i - 20" is "-20 - 9i".  Conjugate laws Here are a few simple rules regarding the complex conjugate and The above laws simply says that the sum of conjugates equals the conjugate of the sum; and similarly, the conjugate of the product equals the product of the conjugates. Consider this example: and we can see that which equals to This confirms the addition conjugate law. The complex root. Now that you are equipped with all the basics of complex numbers, you can tackle the more advanced topic of root finding. Consider the question: Express "w" in the form of "a + ib". That is easy enough. Solve (1) and (2) simultaneously to work out "a" and "b". Observe that if, after solving for "a" and "b", we replace them with -"a" and -"b" respectively, then they would still satisfy the two simultaneous equations above, we can see that (as expected) if "w" = "a" + "ib" satisfies the equation "w"2 = "z", then so will "w" = -("a" + "ib"). With real numbers, we always take the non-negative answer and call the solution formula_50. However, since there is no notion of "greater than" or "less than" with complex numbers, there is no such choice of formula_51. In fact, which square root to take as "the" value of formula_51 depends on the circumstances, and this choice is very important to some calculations. info -- Finding the square root. Finding the root of a real number is a very difficult problem to start with. Most people have no hope of finding a close estimate of formula_53 without the help of a calculator. The modern method of approximating roots involves an easy to understand and ingenius piece of mathematics called the Taylor series expansion. The topic is usually covered in first year university maths as it requires an elementary understanding of an important branch of mathematics called calculus. The Newton-Raphson method of root finding is also used extensively for this purpose. Now consider the problem Express "w" in the form of ""a" + "ib"". Using the methodology developed above we proceed as follows, It turns out that the simultaneous equations (1) &amp; (2) are hard to solve. Actually, there is an easy way to calculate the roots of complex numbers called the De Moivre's theorem, it allows us to calculate the "n"th root of any complex number with ease. But to set the method, we need understand the geometric meaning of a complex number and learn a new way to "represent" a complex number. The complex plane. Complex numbers as ordered pairs. It is worth noting, at this point, that every complex number, "a" + "bi", can be completely specified with exactly two real numbers: the "real part" "a", and the "imaginary part" "b". This is true of "every" complex number; for example, the number 5 has real part 5 and imaginary part 0, while the number 7"i" has real part 0 and imaginary part 7. We can take advantage of this to adopt an alternative scheme for writing complex numbers: we can write them as ordered pairs, in the form "(a, b)" instead of "a+bi". These should look familiar: they are exactly like the ordered pairs we use to represent points in the plane. In fact, we can use them that way; the plane which results is called the "complex plane". We refer to its x axis as the "real axis", and to its y axis as the "imaginary axis". The complex plane. We can see from the above that a single complex number is a point in the complex plane. We can also represent "sets" of complex numbers; these will form "regions" on the plane. For example, the set is a square of edge length 2 centered at the origin (See following diagram). Complex-valued functions. Just as we can make functions which take "real" values and output "real" values, so we can create functions from complex numbers to real numbers, or from complex numbers to complex numbers. These latter functions are often referred to as "complex-valued" functions, because they evaluate to (output) a complex number; it is implicit that their argument (input) is complex as well. Since complex-valued functions map complex numbers to other complex numbers, and we have already seen that complex numbers correspond to points on the complex plane, we can see that a complex-valued function can turn regions on the complex plane into other regions. A simple example: the function formula_58 takes a point in the complex plane and shifts it up by 1. If we apply it to the set of points making up the square above, it will move the entire square up one, so that it "rests" on the x-axis. de Moivre's Theorem. For example, say we have formula_60. We may now write this in polar form, Using de Moivre's theorem, we can deduce, etc. Complex root of unity. The complex roots of unity to the nth degree is the set of solutions to the equation formula_65. Clearly they all have magnitude 1. They form a cyclic group under multiplication. For any given formula_20, there are exactly formula_20 many of them, and they form a regular n-gon in the complex plane over the unit circle. A closed form solution can be given for them, by use of Euler's formula: formula_68 The sum of the formula_20th roots of unity is equal to 0, except for formula_70, where it is equal to 1. The product of the formula_20th roots of unity alternates between -1 and 1. Problem set. \right)^{2i} = 2^ie^{\frac{\pi}{2}}&lt;/math&gt; 

Well done! This is the correct answer. Applying Ohm's law to the top branch: I=V/R =6/4 A =1.5A Applying Ohm's law to the bottom branch: I=V/R =6/(2+2) A =1.5A On to the next question» 

I'm sorry but this is the wrong answer. Remember that each branch has the entire voltage across it. Applying Ohm's law to the top branch: I=V/R =6/4 A =1.5A The lower branch has a combined resistance of 4Ω because the two resistirs are in series. Applying Ohm's law to the bottom branch: I=V/R =6/(2+2) A =1.5A On to the next question» 

I'm sorry but this is the wrong answer. Remember that each branch has the entire voltage across it. Applying Ohm's law to the top branch: I=V/R =6/4 A =1.5A The lower branch has a combined resistance of 4Ω because the two resistirs are in series. Applying Ohm's law to the bottom branch: I=V/R =6/(2+2) A =1.5A On to the next question» 

I'm sorry but this is the wrong answer. Remember that each branch has the entire voltage across it. Applying Ohm's law to the top branch: I=V/R =6/4 A =1.5A The lower branch has a combined resistance of 4Ω because the two resistirs are in series. Applying Ohm's law to the bottom branch: I=V/R =6/(2+2) A =1.5A On to the next question» 

&lt; Back to Questions A plug contains three wires. Each a different colour: Red, Blue, and green/yellow. Only two of the wires in a plug carry current under normal conditions. These are: 

This is the correct answer. The red is the live wire and the blue is neutral. Both wires carry current under normal working conditions. On to the next question» 

I'm sorry but this is the wrong answer. The green/yellow is the Earth wire. The Earth wire does not carry any current under normal operating conditions. It is a safety device. On to the next question» 

I'm sorry but this is the wrong answer. The green/yellow is the Earth wire. The Earth wire does not carry any current under normal operating conditions. It is a safety device. On to the next question» 

A student rubs a polythene strip with fur and suspends it from a clamp stand. She then rubs another polythene strip with fur and brings it up to the first strip. The two strips repel each other. Look at the following four statements: Which two of the above statements are false? 



Personal Pronouns in English. Pronouns are nouns which are used instead of another noun ('pro', in place of 'noun', noun.) There are three categories of pronouns which are divided up into persons: 1st, 2nd, and 3rd. In addition, pronouns can be singular or plural. They are declined like all other nouns. Personal Pronouns in Latin. 1st/2nd Person Pronouns. Table of Personal Pronouns in all of their cases: I, thou, we, ye. Note: ' is the archaic singular of the archaic plural ' - useful for distinguishing "you" (singular) from "you" (plural) Nota Bene: the genitive is used in certain phrases like: For the possessive uses (my sister, your bicycle) sometimes uses the possessive adjectives: 3rd Person Pronouns. Technically, 3rd person pronouns do not exist in Latin as they do in English. However, they do have equivalents. Adjectives modify nouns and take the gender of the noun which they modify. However, adjectives do not necessarily need a substantive present in the sentence to modify. The substantive can be presumed. In this way, '3rd person' pronouns are formed. Example 1. Take the masculine form of the adjective 'ille'. Literally it means 'That (masculine) thing.' However one could take it for simply meaning 'he', depending on the context. Similarly, the pronoun 'iste' means 'that (masc.) thing'. Iste and ille are declined in exactly the same way, but there are a slight difference of meaning between them: 'ille' is often used with proper names for marking dignity or worth and 'iste' conveys a contemptuous sense. Examples: - Annibal, ille inclytus filius Amilcaris (Hannibal, that renowned Hamilcar's son). - Iste servus improbus ante te (this bad slave in front of you). If no substantive is provided assume words like these: 'man', 'woman', 'thing', 'idea', 'concept', 'reason' etc. Let context be your guide. Common Adjectives Used as 3rd Person Pronouns In Latin. Declension of Ille (that). Ille is often used as a kind of pronoun. In situations with multiple phrases or sentences, however, it is syntactically different from is, ea, id (see below). For example: "Canis puero cibum dat. Is laborat in agro." means "The dog gives food to the boy. The dog works in the field". However: "Canis puero cibum dat. Ille laborat in agro." means "The dog gives food to the boy. The boy works in the field". Thus, ille, unlike the other pronouns makes a previous object into the subject (and vice versa). Declension of Is, ea, id: (personal pronouns w/ translations). Like ille, is can be used as a form of a pronoun. Uses of the Relative Pronoun. The relative pronoun takes on the case depending on the function it serves in the relative clause. For example, in the sentence "He sees the man who has a slave," "who" is translated as nominative because it is the subject of the clause "who has a slave." The antecedent (noun to which the pronoun refers) is usually before the relative clause. Declension of hic, haec, hoc (meaning "this"). N.B. Hic as an adverb that means 'here'. N.B. Hic can also be used as a pronoun. 

Im sorry but this is the wrong answer. B is quite possibly true, both rods could be negative because two like charges repel. Hit back and try again. «back 

I'm sorry but this is the wrong answer. C is quite possibly true! two positively charged rods would repel each other. Hit back and try again. «back 

This is the correct answer. Well done! Next question» 

I'm sorry but this is the wrong answer. If both strops are negative they repel. If ther are both positive they also repel. But the question asked which statements are false. Hit back and try again. «back 

I'm sorry but this is the wrong answer. B may very well be true as "like charges repel". Hit back and try again. «back 

I'm sorry but this is the wrong answer. C may well be true. Two positively charged strips would repel. «back 

Acknowlegement. For assistance in preparation of this document, special thanks is extended to Authorization. Resolution Introduction. Introduction Historical Note on Formation of the Constitution. Historical Note on Formation of the Constitution [The Constitution begins with a Preamble. It is followed by Articles, which are divided into Sections, which are again divided into clauses.] Articles. Article Two - Executive Branch. Section 4. - Impeachment. The President, Vice President and all civil Officers of the United States, shall be removed from Office on Impeachment for, and Conviction of, Treason, Bribery, or other high Crimes and Misdemeanors. Article Five - The Amendment Process. The Congress, whenever two thirds of both houses shall deem it necessary, shall propose amendments to this Constitution, or, on the application of the legislatures of two thirds of the several states, shall call a convention for proposing amendments, which, in either case, shall be valid to all intents and purposes, as part of this Constitution, when ratified by the legislatures of three fourths of the several states, or by conventions in three fourths thereof, as the one or the other mode of ratification may be proposed by the Congress; provided that no amendment which may be made prior to the year one thousand eight hundred and eight shall in any manner affect the first and fourth clauses in the ninth section of the first article; and that no state, without its consent, shall be deprived of its equal suffrage in the Senate. Article Seven - Ratificiation Process. The ratification of the conventions of nine states, shall be sufficient for the establishment of this Constitution between the states so ratifying the same. The Amendments. Bill of Rights. Amendment II - Right to Bear Arms (1791). A well regulated militia, being necessary to the security of a free state, the right of the people to keep and bear arms, shall not be infringed. Amendment IV - Search and Seizure (1791). The right of the people to be secure in their persons, houses, papers, and effects, against unreasonable searches and seizures, shall not be violated, and no warrants shall issue, but upon probable cause, supported by oath or affirmation, and particularly describing the place to be searched, and the persons or things to be seized. Amendment VII - Common Law Suits - Jury Trial (1791). In suits at common law, where the value in controversy shall exceed twenty dollars, the right of trial by jury shall be preserved, and no fact tried by a jury, shall be otherwise reexamined in any court of the United States, than according to the rules of the common law. Amendment VIII - Excess Bail or Fines, Cruel and Unusual Punishment (1791). Excessive bail shall not be required, nor excessive fines imposed, nor cruel and unusual punishments inflicted. Amendment IX - Non-Enumerated Rights (1791). The enumeration in the Constitution, of certain rights, shall not be construed to deny or disparage others retained by the people. Amendment X - Rights Reserved to States (1791). The powers not delegated to the United States by the Constitution, nor prohibited by it to the states, are reserved to the states respectively, or to the people. Other Amendments. Amendment XI - Suits Against a State (1795). The judicial power of the United States shall not be construed to extend to any suit in law or equity, commenced or prosecuted against one of the United States by citizens of another state, or by citizens or subjects of any foreign state. Amendment XIII - Abolition of Slavery (1865). Section 1. Neither slavery nor involuntary servitude, except as a punishment for crime whereof the party shall have been duly convicted, shall exist within the United States, or any place subject to their jurisdiction. Section 2. Congress shall have power to enforce this article by appropriate legislation. Amendment XIV - Privileges and Immunities, Due Process, Equal Protection, Apportionment of Representatives, Civil War Disqualification and Debt (1868). Section 1. Section 2. Section 3. No person shall be a Senator or Representative in Congress, or elector of President and Vice President, or hold any office, civil or military, under the United States, or under any state, who, having previously taken an oath, as a member of Congress, or as an officer of the United States, or as a member of any state legislature, or as an executive or judicial officer of any state, to support the Constitution of the United States, shall have engaged in insurrection or rebellion against the same, or given aid or comfort to the enemies thereof. But Congress may by a vote of two-thirds of each House, remove such disability. Section 4. Section 5. The Congress shall have power to enforce, by appropriate legislation, the provisions of this article. Amendment XV - Rights Not to Be Denied on Account of Race (1870). Section 1. The right of citizens of the United States to vote shall not be denied or abridged by the United States or by any state on account of race, color, or previous condition of servitude. Section 2. The Congress shall have power to enforce this article by appropriate legislation. Amendment XVI - Income Tax (1913). The Congress shall have power to lay and collect taxes on incomes, from whatever source derived, without apportionment among the several states, and without regard to any census or enumeration. Amendment XVII - Election of Senators (1913). The Senate of the United States shall be composed of two Senators from each state, elected by the people thereof, for six years; and each Senator shall have one vote. The electors in each state shall have the qualifications requisite for electors of the most numerous branch of the state legislatures. When vacancies happen in the representation of any state in the Senate, the executive authority of such state shall issue writs of election to fill such vacancies: Provided, that the legislature of any state may empower the executive thereof to make temporary appointments until the people fill the vacancies by election as the legislature may direct. This amendment shall not be so construed as to affect the election or term of any Senator chosen before it becomes valid as part of the Constitution. Amendment XVIII - Prohibition (1919). Section 1. After one year from the ratification of this article the manufacture, sale, or transportation of intoxicating liquors within, the importation thereof into, or the exportation thereof from the United States and all territory subject to the jurisdiction thereof for beverage purposes is hereby prohibited. Section 2. The Congress and the several states shall have concurrent power to enforce this article by appropriate legislation. Section 3. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of the several states, as provided in the Constitution, within seven years from the date of the submission hereof to the states by the Congress. Amendment XIX - Women's Right to Vote (1920). The right of citizens of the United States to vote shall not be denied or abridged by the United States or by any state on account of sex. Congress shall have power to enforce this article by appropriate legislation. Amendment XX - Presidential Term and Succession (1933). Section 1. The terms of the President and Vice President shall end at noon on the 20th day of January, and the terms of Senators and Representatives at noon on the 3d day of January, of the years in which such terms would have ended if this article had not been ratified; and the terms of their successors shall then begin. Section 2. The Congress shall assemble at least once in every year, and such meeting shall begin at noon on the 3d day of January, unless they shall by law appoint a different day. Section 3. If, at the time fixed for the beginning of the term of the President, the President elect shall have died, the Vice President elect shall become President. If a President shall not have been chosen before the time fixed for the beginning of his term, or if the President elect shall have failed to qualify, then the Vice President elect shall act as President until a President shall have qualified; and the Congress may by law provide for the case wherein neither a President elect nor a Vice President elect shall have qualified, declaring who shall then act as President, or the manner in which one who is to act shall be selected, and such person shall act accordingly until a President or Vice President shall have qualified. Section 4. The Congress may by law provide for the case of the death of any of the persons from whom the House of Representatives may choose a President whenever the right of choice shall have devolved upon them, and for the case of the death of any of the persons from whom the Senate may choose a Vice President whenever the right of choice shall have devolved upon them. Section 5. Sections 1 and 2 shall take effect on the 15th day of October following the ratification of this article. Section 6. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of three-fourths of the several states within seven years from the date of its submission. Amendment XXI - Repeal of Prohibition (1933). Section 1. The eighteenth article of amendment to the Constitution of the United States is hereby repealed. Section 2. The transportation or importation into any state, territory, or possession of the United States for delivery or use therein of intoxicating liquors, in violation of the laws thereof, is hereby prohibited. Section 3. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by conventions in the several states, as provided in the Constitution, within seven years from the date of the submission hereof to the states by the Congress. Amendment XXII - Two Term Limit on President (1951). Section 1. No person shall be elected to the office of the President more than twice, and no person who has held the office of President, or acted as President, for more than two years of a term to which some other person was elected President shall be elected to the office of the President more than once. But this article shall not apply to any person holding the office of President when this article was proposed by the Congress, and shall not prevent any person who may be holding the office of President, or acting as President, during the term within which this article becomes operative from holding the office of President or acting as President during the remainder of such term. Section 2. This article shall be inoperative unless it shall have been ratified as an amendment to the Constitution by the legislatures of three-fourths of the several states within seven years from the date of its submission to the states by the Congress. Amendment XXIII - Presidential Vote in District of Columbia (1961). Section 1. The District constituting the seat of government of the United States shall appoint in such manner as the Congress may direct: A number of electors of President and Vice President equal to the whole number of Senators and Representatives in Congress to which the District would be entitled if it were a state, but in no event more than the least populous state; they shall be in addition to those appointed by the states, but they shall be considered, for the purposes of the election of President and Vice President, to be electors appointed by a state; and they shall meet in the District and perform such duties as provided by the twelfth article of amendment. Section 2. The Congress shall have power to enforce this article by appropriate legislation. Amendment XXIV - Poll Tax (1964). Section 1. The right of citizens of the United States to vote in any primary or other election for President or Vice President, for electors for President or Vice President, or for Senator or Representative in Congress, shall not be denied or abridged by the United States or any state by reason of failure to pay any poll tax or other tax. Section 2. The Congress shall have power to enforce this article by appropriate legislation. Amendment XXV - Presidential Succession (1967). Section 1. In case of the removal of the President from office or of his death or resignation, the Vice President shall become President. Section 2. Whenever there is a vacancy in the office of the Vice President, the President shall nominate a Vice President who shall take office upon confirmation by a majority vote of both Houses of Congress. Section 3. Whenever the President transmits to the President pro tempore of the Senate and the Speaker of the House of Representatives his written declaration that he is unable to discharge the powers and duties of his office, and until he transmits to them a written declaration to the contrary, such powers and duties shall be discharged by the Vice President as Acting President. Section 4. Whenever the Vice President and a majority of either the principal officers of the executive departments or of such other body as Congress may by law provide, transmit to the President pro tempore of the Senate and the Speaker of the House of Representatives their written declaration that the President is unable to discharge the powers and duties of his office, the Vice President shall immediately assume the powers and duties of the office as Acting President. Thereafter, when the President transmits to the President pro tempore of the Senate and the Speaker of the House of Representatives his written declaration that no inability exists, he shall resume the powers and duties of his office unless the Vice President and a majority of either the principal officers of the executive department or of such other body as Congress may by law provide, transmit within four days to the President pro tempore of the Senate and the Speaker of the House of Representatives their written declaration that the President is unable to discharge the powers and duties of his office. Thereupon Congress shall decide the issue, assembling within forty-eight hours for that purpose if not in session. If the Congress, within twenty-one days after receipt of the latter written declaration, or, if Congress is not in session, within twenty-one days after Congress is required to assemble, determines by two-thirds vote of both Houses that the President is unable to discharge the powers and duties of his office, the Vice President shall continue to discharge the same as Acting President; otherwise, the President shall resume the powers and duties of his office. Amendment XXVI - Right to Vote at Age 18 (1971). Section 1. The right of citizens of the United States, who are 18 years of age or older, to vote, shall not be denied or abridged by the United States or any state on account of age. Section 2. The Congress shall have the power to enforce this article by appropriate legislation. Amendment XXVII - Compensation of Members of Congress (1992). No law, varying the compensation for the services of the Senators and Representatives, shall take effect, until an election of Representatives shall have intervened. References.  

The Election of 1840. President Martin Van Buren was blamed for the Panic of 1837, but felt that he deserved to be reelected in 1840. Van Buren was a Democrat from New York who had continued the policies of Andrew Jackson. To oppose him the Whig Party joined to bring in a hero of the Indian wars, William Henry Harrison, "Old Tippecanoe." The ticket was balanced by the Vice Presidential candidate, a Southerner named John Tyler. The Harrison campaign was thoroughly managed. The campaign song, "Tippecanoe and Tyler Too," was headed with an image of the log cabin where Harrison had supposedly grown up. Paid staffers went to frontier towns rolling a huge canvas ball, inscribed, "Keep it rolling for Tippecanoe and Tyler Too." (The American idiom "Keep the ball rolling" comes from this usage.) The ball would stop in front of a local tavern, then a common meeting place of the community. There they would stage a rally, typically with some free cider. Another sign of Harrison's plain man status was his title as "the hard cider candidate." There was little discussion of the issues. For its part, Van Buren's campaign called Harrison a provincial, out-of-touch old man. (The latter was then sixty-eight years old, a rare age in those days.) Harrison won, and gave an hours-long, polished inaugural speech to prove his sophistication. Three weeks afterward, he came down with a cold which turned into pneumonia. He died in April of 1841, and John Tyler was sworn in as President. Thus the Whig Party, predominately Northern and ambivalent about slavery, elected a Virginian advocate of slavery and opponent of the American System. This was a startling omen for those like Clay who believed in American unity. John Tyler Presidency. Tyler's dislike of Jackson had moved him to change his party from Democrat to Whig. His government marked the only Whig presidency. His supporters included formerly anti-Jackson Democrats and National Republicans. He supported states' rights; so when many of the Whig bills came to him, they were never voted in. In fact, Tyler vetoed the entire Whig congressional agenda. The Whigs saw this as a party leader turning on his own party. He was officially expelled from the Whig party in 1841. The Tyler presidency threw the Whig party into disarray. Because of divisions between the two factions in the party, the Whigs could not agree on one goal. Much of the public did not take Tyler's presidency seriously. They saw his lack of appeal in Congress and the embarrassing resignations of all of but one of Harrison's cabinet appointees in a single month. Yet Tyler's administration helped polarize the two parties. When he appointed John C. Calhoun, a staunch pro-slavery Democrat, as his Secretary of State, he confirmed a growing feeling that Democrats were the party of the South and Whigs the party of the North. In the election of 1844, Whigs voted by sectional ties. Because of these weakening divisions within the party, the Democratic candidate, James Polk, won. After one term, the Whigs were out of power. Manifest Destiny. Many Western European-descended "White" Americans supported anti-Native American policies. The theme of conquest over the Indian was seen as early as John Filson's story of Daniel Boone in 1784. In the Nineteenth Century this was joined to the conviction that the United States was destined to take over the whole continent of North America, the process of Manifest Destiny articulated by John O' Sullivan in 1845.[source needed] America carried the Bible, civilization, and democracy: the Indian had none of these. Many European descendants believed other ethnic groups, including those people imported as slaves from Africa and their descendants, were childlike, stupid, and feckless. It was the duty of so-called superior groups to meet these inferior groups and to dominate them. So-called inferior ethnic groups could not advance technologically or spiritually. The idea of Manifest Destiny resulted in the murders and dislocation of millions of people. The Cherokee had been converted to Christianity, they were by-and-large peaceful, and they were using a self-invented alphabet to print newspapers. But their deportation, the "Trail of Tears," was justified by Manifest Destiny. The conviction was behind the Louisiana Purchase, the final shaking of French colonialism in what would become the Continental United States. It was behind the defeat of Spanish and Mexicans in a succession of skirmishes and wars. It helped send out pro- and anti-slavery factions across new areas, and still later brought about legislation such as the Homestead Act. Amistad Case. In February of 1839, Portuguese slave hunters abducted a large group of Africans from Sierra Leone and shipped them to Havana, Cuba, a center for the slave trade. This abduction violated all of the treaties then in existence. Fifty-three Africans were purchased by two Spanish planters and put aboard the Cuban schooner Amistad for shipment to a Caribbean plantation. On July 1, 1839, the Africans seized the ship, killed the captain and the cook, and ordered the planters to sail to Africa. On August 24, 1839, the Amistad was seized off Long Island, NY, by the U.S. brig Washington. The planters were freed and the Africans were imprisoned in New Haven, CT, on charges of murder. Although the murder charges were dismissed, the Africans continued to be held in confinement as the focus of the case turned to salvage claims and property rights. President Van Buren was in favor of extraditing the Africans to Cuba. However, abolitionists in the North opposed extradition and raised money to defend the Africans. Claims to the Africans by the planters, the government of Spain, and the captain of the brig led the case to trial in the Federal District Court in Connecticut. The court ruled that the case fell within Federal jurisdiction and that the claims to the Africans as property were not legitimate because they were illegally held as slaves. The case went to the Supreme Court in January 1841, and former President John Quincy Adams argued the defendants' case. Adams defended the right of the accused to fight to regain their freedom. The Supreme Court decided in favor of the Africans, and 35 of them were returned to their homeland. The others died at sea or in prison while awaiting trial. The result, widely publicized court cases in the United States helped the abolitionist movement. Technology. The canals had been a radical innovation. But they had their limitations. They could only overcome mountains with complicated, overland bypasses, and in winter they froze, stopping traffic completely. But an answer was found in the steam-driven, coal-powered engine. The steamboat was already bringing cotton and people up the rivers, erasing an age-old transportation problem. The development of railroad engines made travel and manufacture possible even in winter. It made the expensive canal obsolete: wherever you could run a rail, you could have a town. And the coal-fired, steam-powered engine could bring manufacturing to places without great rivers. The prosperity of the New England mill towns could be replicated elsewhere. Coal and its byproducts became a major industry in America. (In the 1850s some German cities became known for creating coal-based dyes to make bold-colored fabrics.) Iron works and glass plants built large furnaces, fueled by coke, a coal derivative. They were contained by huge buildings. Steam-boats burned coke. So did steam-driven works. Smoke and smut from industry and household coal fires poured into city air. In his 1842 tour of Pittsburgh, Charles Dickens looked at the haze and fire and called it "Hell with the lid off." Canals, railroads, and the teletype system tied the country together in a way thought impossible in 1790. They increased the market for goods, and thus the demand. The Second Industrial Revolution produced faster ways of satisfying that demand. In 1855 Henry Bessemer patented a furnace which could turn iron into steel, in high quantity. Iron workers, "puddlers," had worked slowly and regularly, had been paid a high wage, and had been considered craftsmen. The new steelworkers did not need that skill. They could be paid more cheaply. In other industries, faster processes of work either made a mockery of the apprenticeship system or eliminated it altogether. Manufacturers faced the same situation of the New England cloth makers a generation before, and solved it another way. To find fresh, fast, cheap labor, you could often hire children. You didn't have to pay them as much, and they didn't complain. Where children used to work on the farm, they now worked in groups in factories for higher wages. These children were, at the best, expected to work long and hard. (It was cheaper to run big machines in shifts than to have them idle all night.) Boys and girls worked naked in the coal mines; boys got burned in the glassworks; boys got maimed or killed running heavy machinery. None of them went to school, so even the ones who survived to adulthood were unfit for jobs when they came of age. The ideals of Thomas Jefferson were dead. Instead of craftsmen and farmers living by their own hands, the cities were being filled by people who owned little or nothing, getting on by the wages paid by often indifferent employers. This was true in New England and the Middle States, though not in the South (except for the Tregar Ironworks in Richmond, Virginia, itself partially manned by slaves). Politicians such as John C. Calhoun jeered at Northern "wage slaves," and dreamed of a South with the technology and government of Sparta. Compromise of 1850. The Compromise of 1850 was an intricate package of five bills passed in September 1850. It defused a four-year confrontation between the slave states of the South and the free states of the North that arose following the Mexican-American War. The compromise, drafted by Whig Henry Clay and brokered by Democrat Stephen Douglas, quieted sectional conflict for four years. The calm was greeted with relief, although each side disliked specific provisions. Texas surrendered its claim to New Mexico, but received debt relief and the Texas Panhandle, and retained the control over El Paso that it had established earlier in 1850. The South avoided the humiliating Wilmot Proviso, but did not receive desired Pacific territory in Southern California or a guarantee of slavery south of a territorial compromise line like the Missouri Compromise Line or the 35th parallel north. As compensation, the South received the possibility of slave states by popular sovereignty in the new New Mexico Territory and Utah Territory, which, however, were unsuited to plantation agriculture and populated by non-Southerners; a stronger Fugitive Slave Act, which in practice outraged Northern public opinion; and preservation of slavery in the national capital, although the slave trade was banned there except in the portion of the District of Columbia that adjoined Virginia. The Compromise became possible after the sudden death of President Zachary Taylor, who, although a slave owner himself, tried to implement the Northern policy of excluding slavery from the Southwest. Whig leader Henry Clay designed a compromise, which failed to pass in early 1850. In the next session of Congress, Democratic Senator Stephen Douglas of Illinois narrowly passed a slightly modified package over opposition by extremists on both sides, including Senator John C. Calhoun of South Carolina. Texas and Mexico. Mexico won independence from Spain in 1821. Weakened by more than a decade of struggle, the new Republic of Mexico attempted to attract settlers from the United States to the then-sparsely populated Mexican state of Coahuila y Texas. The first white settlers were 200 families led by Stephen F. Austin as a part of a business venture started by Austin's father. Despite nominal attempts to ensure that immigrants would be double penetrated with Mexican cultural values -- by requiring, for example, acceptance of Catholicism and a ban on slave holding -- Mexico's immigration policy led to the whites, rather than Mexicans, becoming the demographic majority in Texas by the 1830's, their beliefs and American values intact. Due to past US actions in Texas, Mexico feared that white Americans would convince the United States to annex Texas and Mexico. In April 1830, Mexico issued a proclamation that people from the United States could no longer enter Texas. Mexico also would start to place custom duties on goods from the United States. In October 1835, white colonists in Texas revolted against Mexico by attacking a Mexican fort at Goliad, defeating the Mexican garrison. At about the same time, the Mexican president, Antonio López de Santa Anna, provoked a constitutional crisis that was among the causes of the revolt in Texas, as well as a rebellion in the southern Mexican province of Yucután. An official declaration of Texas independence was signed at Goliad that December. The next March, the declaration was officially enacted at the Texan capital of Washington-on-the-Brazos, creating the Republic of Texas. A few days before the enactment of the declaration, a Mexican force led by General Antonio López de Santa Anna laid siege to the Alamo, a mission in present day San Antonio. Vastly outnumbered, fewer than 200 Texans at San Antonio de Béxa, renamed the Alamo, held out for 12 days, until the final attack at dawn on March 6, 1836. Santa Anna, as he had promised during the siege, killed the few prisoners taken in the capture. Though the Alamo had been garrisoned in contravention of orders from Sam Houston, who had been placed in charge of Texan armed forces, the delay their defense forced on the Mexican army allowed the Texan government some crucial time to organize. The next month saw the battle of San Jacinto, the final battle of the Texas Revolution. A force of 800 led by Sam Houston, empowered by their rallying war cry of "Remember the Alamo!", defeated Santa Anna's force of 1600 as they camped beside the sluggish creek for which the 20-minute-long battle is named. Santa Anna himself was captured and the next day signed the Treaties of Velasco, which ended Mexico-Texas hostilities. After the fighting had ended, Texas asked to be admitted to the Union, but Texas's request forced Congress to an impasse. One of the most significant problems with the annexation of Texas was slavery. Despite Mexican attempts to exclude the practice, a number of white-Texans held slaves, and the new Republic of Texas recognized the practice as legitimate. In the United States, The Missouri Compromise of 1818 provided for an equality in the numbers of slave and non-slave states in the US, and to allow Texas to join would upset that power balance. For about ten years, the issue was unresolved, until President James Polk agreed to support the annexation of Texas. In 1845, Texas formally voted to join the US. The Mexicans, however, who had never formally recognized Texas's independence, resented this decision. The southern boundary with Texas had never officially been settled and when the United States moved federal troops into this disputed territory, war broke out (assisted by raids carried out across the border by both sides). In the Mexican-American War, as this was called, the US quickly defeated the Mexican Army by 1848. The peace settlement, called the Treaty of Guadalupe Hidalgo, ceded one-third of Mexico's territory to the United States. In addition to Texas, with the border fixed at the Rio Grande River, the United States acquired land that would become the present-day states of New Mexico, California, Arizona, Colorado, Nevada, Utah, and parts of Colorado and Wyoming; the US paid Mexico $15 million. However, the new territories posed even more problems relating to slavery: the balance between slave and non-slave states seemed threatened again. Oregon. In 1824 and 1825 Russia gave up its claim to Oregon. The U.S. and Canada jointly made an agreement for occupation. However, disputes surfaced over the northwestern boundary of the US and the southwestern boundary of Canada. The US claimed that it owned land south of Alaska, while the British claimed that the boundary was drawn at present-day Oregon. President Polk, who had initiated the dispute, gave Great Britain an ultimatum -- negotiate or go to war. On June 15, 1846, Britain agreed to give up the land south of the 49th parallel, while keeping Vancouver Island and navigation rights to the Columbia River. Polk agreed. Comparing this incident to the president's aggressiveness toward Mexico, several individuals [whom?] concluded that Polk favored the causes of the South over those of the North. Oregon Trail Sometimes Native Americans and white settlers met in peace. During the twenty years after 1840, around 250,000 to 500,000 people walked the Oregon Trail across most of the continent on foot, with the trek taking an average of seven months. Many of these settlers were armed in preparation for Native attack, but the majority of the encounters were peaceful. Most of the starting points were along the Missouri River, including Independence, St. Joseph, and Westport, Missouri. Many settlers set out on organized wagon trains, while others went on their own. Settlers timed their departures so they would arrive after spring, allowing their livestock days of pasture at the end, and yet early enough to not travel during the harsh winter. Walking beside their wagons, settlers would usually cover fifteen miles a day. Men, women and children sometimes endured weather ranging from extreme heat to frozen winter in their 2,000 mile journey West. If a traveler became ill, he or she would have no doctor and no aid apart from fellow travelers. Only the strong finished the trail. Although most interactions between Native Americans and settlers were undertaken in good faith, sometimes things went bad. Eventually hostile relations would escalate into full blown war and many years of bloodshed. California. California Territory. When war broke out between the United States and Mexico in 1845, a few white settlers in the Sacramento Valley in the Mexican state of California seized the opportunity to advance white business interests by declaring independence from Mexico despite the wishes of many Mexicans and natives present in California. Before the arrival of Europeans, scholars place the population of California at 10 million natives. The sparsely populated Bear Flag Republic, as the new nation was called, quickly asked the US for protection from Mexico, allowing US military operations in the new Republic's territory. As skirmishes occurred in California, Mexicans suffered many abuses at the hands of the new white government. When the war ended, the California territory and a large surrounding territory were ceded by Mexico to the US in exchange for $15 million. The territory included what would become present day California, Nevada, Utah, most of New Mexico, Arizona, and Colorado and a small part of Wyoming. The continental US was nearly complete. The final piece would come in 1853, when southern Arizona and New Mexico were bought from Mexico for $10 million. The land from the purchase, known as the Gadsden Purchase, was well suited for building a southern transcontinental railroad. California Gold Rush. In 1848 gold was found at the mill of John Sutter, who lived in the foothills of the Sierra Nevada mountain range, 40 miles east of Sacramento. Word of the gold on the American River (the river on which Sutter's mill was located on) spread, and hordes of people rushed into California to mine gold. The rush peaked in 1849, and those who came during that year were known as "forty-niners." The population of the northern California city of San Francisco exploded as a result of the immigration to the region. Many immigrants that joined the Gold Rush did not find opportunity but rather discrimination at the hands of white prospectors and newly changed government. One of these, Joaquin Murrieta, known as the Mexican Robin Hood, had become a bandit and hero of those still loyal to Mexico. As a reaction the Governor of California, John Bigler, formed the California Rangers. This group went after and allegedly found Murrieta and his companions. They cut off his head, which was later put on display. Many still doubt whether the person the California Rangers decapitated was actually Murrieta or some other poor soul. Be that as it may, the memory of Murrieta is still much loved and respected by Mexican Americans today. Apart from being gained by a handful of very lucky prospectors, a great deal of the wealth generated by the Gold Rush belonged to those who owned businesses that were relevant to gold mining. For example, Levi Strauss, a German Jew, invented denim pants for prospectors when he observed that normal pants could not withstand the strenuous activities of mining. Strauss eventually became a millionaire, and the Levi's brand still is recognized today. Mormonism. The Birth of the Latter Day Saints. One continuation of the Second Great Revival is seen in the birth of an American faith, Mormonism, or The Church of God of the Latter Day Saints. Joseph Smith, a resident of New York State, said that he had found golden plates. The documents which he supposedly translated from these plates revealed what he said was a restoration of the faith which Jesus and the apostles had known, a new American-based order. In 1830 he organized what he designated The Church of Christ, or the Church of the Latter Day Saints. This body spread through conversion, its truth seen in its organization and its prosperity. It had several divergences from existing United States law, including the doctrine that men might have more than one wife. After Smith had been arrested in 1844 in Illinois, charged by civil officials with starting a riot and with treason, a crowd broke into the jail where he was being held and murdered him. The Great Mormon Exodus. Yet the Latter Day Saints persevered. Smith's successor was another prophet, Brigham Young. Continued conflict between the U.S. Government, most signally the state of Illinois, and the Mormons led to the decision to leave the States and go to a less-settled place. The territory of Utah, obtained through the wars with Mexico, certainly counted as less-settled: a vast alkali desert, punctuated by grotesque mountains, and sparsely peopled by Spanish-speaking settlements and Indian tribes. The Mormons began sending out a few pioneers for the new territory as early as 1846. In the two decades afterward, while conversion and population growth further increased the Mormons, about 70,000 people made the trek through difficult conditions. In some cases the migrants walked on foot through hostile landscapes, carrying all their goods with them in handcarts they pulled themselves. When they reached Utah, they formed tightly-organized, top-down structures driven by doctrine and individual discipline. The settlers diverted mountain streams to their fields. In places which had been waste, they created fertile farms and productive vegetable gardens. Continuing Skirmishes. Yet even in this new land, conflicts continued between Mormons and the U.S. government. In the spring of 1857, President James Buchanan appointed a non-Mormon, Alfred Cumming, as governor of the Utah Territory, replacing Brigham Young, and dispatched troops to enforce the order. The Mormons prepared to defend themselves and their property; Young declared martial law and issued an order on Sept. 15, 1857, forbidding the entry of U.S. troops into Utah. The order was disregarded, and throughout the winter sporadic raids were conducted by the Mormon militia against the encamped U.S. army. Buchanan dispatched (Apr., 1858) representatives to work out a settlement, and on June 26, the army entered Salt Lake City, Cumming was installed as governor, and peace was restored. In 1890, the president of the Mormon Church, Wilford Woodruff, ruled that there would be no more plural marriages. Other distinctive practices which had become Church practice under Joseph Smith continued, including theocratic rule and declaration of people as gods. By 1896, when Utah became one of the United States, it was both Mormon and as American as Massachusetts or New York State. Public Schools and Education. The Board of Education in Massachusetts was established in 1837, making it the oldest state board in the United States. Its responsibilities were and are to interpret and implement laws for public education in the Commonwealth of Massachusetts. Public education in the Commonwealth was organized by the regulations adopted by the Board of Education, which were good faith interpretations of Massachusetts and federal law. The Board of Education was also responsible for granting and renewing school applications, developing and implementing the Massachusetts Comprehensive Assessment System, submitting yearly budget proposals for public education to the Massachusetts General Court, setting standards for teachers, certifying teachers, principals, and superintendents, and monitoring achievements of districts in the Commonwealth. There was a movement for reform in public education. The leader of this movement was Horace Mann, a Massachusetts lawyer and reformer. He supported free, tax-supported education to replace church schools and the private schools set up by untrained, young men. Mann proposed universal education, which would help Americanize immigrants. During Mann’s tenure as secretary of the Massachusetts Board of Education from 1837 to 1848, Massachusetts led the common school movement brought training for teachers, lengthened school years and raised the teachers pay to attract people to that profession. In this period education began being extended more and more to women. Early elementary schools had separate rooms for boy students and girl students. Now some elementary classes became co-educational, and women began to be hired as teachers. The first woman's college, Mt. Holyoke, was founded in South Hadley, Massachusetts. It was created by Mary Lyon, and intended as a liberal arts college. "Bleeding Kansas". There was never much doubt that the settlers of Nebraska would, in the face of popular sovereignty, choose to bar slavery. Kansas, however, was another matter. Abolitionist and pro-slavery groups tried to rush settlers to Kansas in hopes of swinging the vote in the group's own direction. Eventually, both a free-state and a slave-state government were functioning in Kansas - both illegal. Violence was abundant. In May 1856, a pro-slavery mob ransacked the chiefly abolitionist town of Lawrence, demolishing private property of the anti-slavery governor, burning printing presses, and destroying a hotel. Two days later, in retaliation, Abolitionist John Brown and his sons went to the pro-slavery town of Pottawatomie Creek and hacked five men to death in front of their families. This set off a guerilla war in Kansas that lasted through most of 1856. Violence over the issue of Kansas was even seen in the Senate. Massachusetts Senator Charles Sumner accused South Carolinian Andrew Butler of having "chosen a mistress to whom he has made his vows - Slavery." Upon hearing these words, Butler's nephew, Representative Preston Brooks, walked onto the Senate floor and proceeded to cane Sumner in the head. Sumner suffered so much damage from the attack that he could not return to the Senate for over three years. Brooks was expelled by the House. Cheered on by southern supporters (many of whom sent Brooks new canes, to show approval of his actions), came back after a resounding reelection. After much controversy and extra legislation, Kansas found itself firmly abolitionist by 1858. Dred Scott v. Stanford - 1857. Dred Scott was an African-American slave who first sued for his freedom in 1846. His case stated that he and his wife Harriet had been transported to both the state of Illinois and Minnesota territory. Laws in both places made slavery illegal. Dred and Harriet began with two separate cases, one for each of them. Slaves were not allowed legal marriage, but the two considered themselves married, and wanted to protect their two teenage daughters. As Dred became ill, the two merged their suit. At first it was the rule that state courts could decide if African-Americans in their jurisdictions were slave or free. After many years and much hesitation, the Supreme Court agreed to hear the case. The United States Supreme Court ruled 7-2 in favor of the slave master, citing precedent that found that neither Dred nor his wife could claim citizenship. As they were not citizens, they did not have a claim in Federal Court. The majority argument cited the Missouri Compromise of 1854 to state that a temporary residence outside of Missouri did not immediately emancipate them, since the owner would be unfairly deprived of property. Ostend Manifesto. Southern slave owners had a special interest in Spanish-held Cuba. Slavery existed on the island, but a recent rebellion in Haiti had spurred some Spanish officials to consider emancipation. The Southerners did not want freed slaves so close to their shores, and other observers thought Manifest Destiny should be extended to Cuba. In 1854 three American diplomats, Pierre Soulé (the minister to Spain), James Buchanan (the minister to Great Britain), and John Y. Mason (the minister to France), met in Ostend, Belgium. They held in common the same views as many Southern Democrats. The diplomats together issued a warning to Spain that it must sell Cuba to the United States or risk having it taken by force. This statement had not been authorized by the Franklin Pierce administration and was immediately repudiated. Reaction, both at home and abroad, was extremely negative. Women's History of the Period. Declaration of Sentiments. 1848 marked the year of the Declaration of Sentiments; it was a document written as a plea for the end of discrimination against women in all spheres of Society. Main credit is given to Elizabeth Cady Stanton for writing the document. The document was presented at the first women's rights convention held at Seneca Falls, New York. Though the convention was attended by 300 women and men, only 100 of them actually signed the document which included; 68 women and 32 men. Elizabeth Blackwell. In 1849 Elizabeth Blackwell became the first woman to receive a medical degree. She attended Geneva College in New York and graduated on January 23, 1849. Even though she had her medical degree she was still banned from practicing in most hospitals. She then relocated to Paris, France and continued her training as a midwife instead of a physician. While in Paris she contracted an eye infection from a small baby that forced her to lose her right eye. It was replaced by a glass eye which ended her medical career. Missouri v. Celia. This murder trial took place in Calloway County, Missouri beginning October 9, 1855. It involved a slave woman named Celia and her master Robert Newsome. After being purchased at the age of 14 in 1850 Celia bore two of her masters children. Soon after becoming intimate with another slave while still being sought after by her master Celia became pregnant. On June 23, 1855, feeling unwell from the pregnancy, Celia pleaded with her master to let her rest; when Newsome ignored her pleas she struck him twice in the head with a heavy stick. She then spent the night burning his corpse in her fireplace and grinding the smaller bones into pieces with a rock. Although Missouri statutes forbade anyone "to take any woman unlawfully against her will and by force, menace or duress, compel her to be defiled," the judge residing over the case instructed the jury that Celia, being enslaved, did not fall within the meaning of "any woman" thus since the "sexual abuser" was her master the murder was not justified on the claim of self-defense. Celia was found guilty of the crime on October 10, 1855 and was sentenced to be hanged. The case still remains significant in history because it graphically illustrates the dreadful truth that enslaved women had absolutely no recourse when it came to being raped by their masters. Panic of 1857. The Panic of 1857 introduced the United States, at least in a small way, to the intricate dealings of the worldwide economy. On the same day that the Central America wrecked, Cincinnati's Ohio Life Insurance and Trust Company ceased operation thanks to embezzlement. News of the twin disasters spread quickly, in part because of the telegraph now becoming common. Investors, including British investors, began to withdraw money from Wall Street in massive numbers. Bank failures increased, mostly in the industrial Northeast and New England states, while the West and South, still more dependent on agriculture, seemed to weather the storm better. There were many underlying causes for the Panic of 1857, and by the time the twin disasters occurred the United States was well on its way into the economic downturn. For 3 years the Crimean War had involved European and Asian countries which increase foreign dependence on American agriculture. The return of the men and land to agricultural production meant an abundance of crops in 1857 which led to falling prices for farm goods. Land speculation, too, had become rampant throughout the United States. This led to an unsustainable expansion of the railroads. As investment money dried up, the land speculation collapsed, as did many of the railroads shortly thereafter. Attempts were made by the federal government to remedy the situation. A bank holiday was declared in October, 1857 and Secretary of the Treasury Howell Cobb recommended the government selling revenue bonds and reducing the tariff (Tariff of 1857). By 1859 the country was slowly pulling out of the downturn, but the effect lasted until the opening shots of the Civil War. Rebellion at Harper's Ferry, Virginia. John Brown. John Brown had been born in Connecticut on May 19, 1800. He grew up in Ohio, where his father worked as a tanner and a minister near Oberlin, Ohio. His father preached abolitionism, and John Brown learned from him. He married twice, his first wife dying while giving birth to their seventh child. He would ultimately father twenty children, eleven of them surviving to adulthood. He started several failed business ventures and land deals in Ohio and Massachusetts. For a while he settled in a community of both black and white settlers. He lived there peacefully until the mid-1850s. Then two of his sons who had moved to Kansas asked him for guns to defend themselves against Missouri Border Ruffians. After two failed defense efforts, Brown left the Kansas area to avoid prosecution for the Pottawatomie massacre. He was already gaining some mention in the press for his efforts. He moved back East and decided to plan a way to destroy slavery in America forever. Brown's Raid On Harper's Ferry. After the troubles in Kansas Brown decided on a plan. A lightly-defended armory in Harper's Ferry, Northern Virginia, contained 100,000 muskets and rifles. An attacker would need some monetary investment to obtain a battalion of men, a similar number of rifles, and a thousand pikes. With the weapons seized at the armory, Brown planned on arming sympathizers and slaves freed by his personal army as it swept through the South. Harper's Ferry had no plantations, and Brown expected no resistance from the local townspeople. On October 16, 1859, Brown carried out his raid, which he planned as the beginning of his revolution. However, instead of his battalion of 450 men, he went in with a group of twenty, including two of his sons. They easily overtook the single nightwatchman and killed several townspeople on the way in, including a free African American man who discovered their plot. Brown had also underestimated the resolve of the local townspeople, who formed a militia and surrounded Brown and his raiders in the armory. After a siege of two days, the U.S. Army sent in a detachment of Marines from Washington, D.C., the closest available contingent. The marines, led by Robert E. Lee, stormed the armory, and in a three minute battle ten of Brown's men were killed. Brown and six others were taken alive and imprisoned for a swift trial. Brown and five of his raiders were hung before the end of the year. Three others were killed in early 1860. Public Reaction. News of the rebellion spread rapidly around the country by telegraph and newspaper, though opinions differed about what it meant. "The Charleston Mercury" of November seven, 1859, represents one Southern view: "With five millions of negroes turned loose in the South, what would be the state of society? It would be worse than the 'Reign of Terror'. The day of compromise is passed." The reaction was most mixed and vigorous among those who called themselves Christians. Abraham Lincoln typified the response of many others when he said that, though Brown "agreed with us in thinking slavery wrong," "[t]hat cannot excuse violence, bloodshed, and treason." The Unitarian William Lloyd Garrison, already having swerved from his previously-held pacifism, on the day of Brown's death preached a sermon commemorating him. "Whenever there is a contest between the oppressed and the oppressor, -– the weapons being equal between the parties, –- God knows that my heart must be with the oppressed, and always against the oppressor. Therefore, whenever commenced, I cannot but wish success to all slave insurrections." The Congregationalist Minister Henry Ward Beecher likewise supported Brown from the pulpit. In short, a faction of White people of faith both conservative and liberal began saying that only by violence could slavery be struck from America. Election of 1860. The new-born Republican party supported Northern industry, as the Whigs had done. It also promised a tariff for the protection of industry and pledged the enactment of a law granting free homesteads to settlers who would help in the opening of the West. But by 1860, it had become the party of abolition. Many Republicans believed that Lincoln's election would prevent any further spread of slavery. It selected Abraham Lincoln of Illinois as its presidential candidate, and Hannibal Hamlin of Maine as its vice-presidential candidate. The Democratic party split in two. The main party or the Northern Democrats could not immediately decide on a candidate, and after several votes, their nominating convention was postponed when the Southern delegates walked out. When it eventually resumed, the party decided on Stephen Douglas of Illinois as its candidate. The first vice-presidential candidate, Benjamin Fitzpatrick, dropped his name from consideration once his home state of Alabama seceded from the Union. His replacement was Herschel Johnson of Georgia. The Southern delegates from the Democratic party selected their own candidate to run for president. John C. Breckenridge of Kentucky with Joseph Lane of Oregon as their vice-presidential candidate. Former Whigs and Southern Republicans who supported the Union in the slavery issue formed the Constitutional Union party. Tennessee senator John Bell was chosen as the Constitutional Union party presidential candidate, over former Texas governor Sam Houston. Harvard President Edward Everett was chosen as the vice-presidential candidate. Abraham Lincoln won the election with only forty percent of the vote. But with the electorate split four ways, it led to a landslide victory in the Electoral College. Lincoln garnered 180 electoral votes without being listed on any of the ballots of any of the future secessionist states in the deep South (except for Virginia, where he received 1.1% of the vote). Stephen Douglas won just under 30% of the popular vote, but only carried 2 states for a total of 12 electoral votes. John Breckenridge carried every state in the deep South, Maryland, and Delaware, for a total of 72 electoral votes. Bell carried the border slave states of Kentucky, Virginia, and Tennessee, for a total of 39 electoral votes. Except for the split decision in the presidential election of 1824, no President in US History has won with a smaller percentage of the popular vote. Lincoln's election ensured South Carolina's secession, along with Southern belief that they now no longer had a political voice in Washington. Other Southern states followed suit. They claimed that they were no longer bound by the Union, because the Northern states had in effect broken a constitutional contract by not honoring the South's right to own slaves as property. Questions For Review. 1. What affect did proslavery sentiments have on the Mexican War? 2. Reconstruct Brown's raid in terms of what happened first, second, and so on. 

Electromagnets are used inside all sorts of devices. A very simple device is the electromagnet that is used in scrap yards to pick up cars. It can then drop them once the electric current is turned off. This page looks at some more complicated devices that use electromagnets. What are Electromagnets? Electromagnets are a type of magnet that you can make using a piece of iron, a wire (that you have to coil around the iron) and a battery. You have to switch on the battery, which should be connected to the wire. The iron will attract other metals, just like a magnet. If you leave the wire attached to the battery, it will begin to get hot. This will also work with a non magnetic core, but the electromagnet will be weaker. The more coils you have, the stronger the electromagnet will be. You can use electromagnets for lots of things: loudspeakers, relays and circuit-breakers. They are also used in car scrap-yards to pick up magnetic waste (e.g. Old cars),cand with any form of magnetically recordable media: VCR tapes, computer hard drives, cassette tapes, and credit card stripes. How Do They Work? Every wire has a magnetic field around it, but this only happens if there is current flowing through it. Every current produces a magnetic field. This effect is used in electromagnets. The field which is produced by the electromagnet is similar to the field around a bar magnet. If it has no iron core it will still produce a field, only it won’t be so strong because the iron makes it stronger. The iron core becomes magnetised when the current is switched on and it loses its magnetism when it is switched off. A steel core however, will always keep its magnetism. The electric bell. Step by Step Explanation. Once the battery is connected a current flows in the wire loops around the U shaped soft iron core. This turns the core into an electromagnet. The electromagnet attracts the armature which is also made of soft iron. As the armature moves towards the electromagnet, it causes the hammer to strike the bell. At the same time it breaks the circuit. The break in the circuit causes the current to stop flowing. The soft iron core loses its magnetic field and releases the armature which springs back to its original position. The contact screw touches the springy metal and completes the circuit. The whole cycle is repeated over and over until the battery is disconnected or runs down. Hope you will find this information helpful! The Loud Speaker. A speaker is made of a paper cone attached to a coil which acts as a small electromagnet. The coil is fitted over a permanent magnet, and as the current flows through the coil it is either attracted or repelled depending on which direction the current flows. A typical signal from an amplifier will be a varying current, and so the cone will vibrate back and forth at the same frequency as the current. As the cone vibrates, it sets up pressure waves in the air, which we hear as sound. 

This book, US Constitution and Government, was initiated and written by , with slight contributions by and . Although the content of the book is comprehensive, it still needs to be edited and revised. The original version of the book was finished on September 30, 2003, making it the first "finished" Wikibook on the site. 

Causes of the Civil War. There are several fundamental causes of the civil war, most of which had related to the south's use of slavery. These include the election of Abraham Lincoln without a single southern electoral college vote. The rise of the Republican party which was opposed to the westward expansion of slavery. The south wanted to protect the rights of their states to determine how they could treat slaves free of federal interference. The northern and southern economies were vastly different, mainly as a result of the south's use of slavery compared to the north's use of free labor which encouraged industrialization. Dixie's Constitution. By the end of March, 1861, the Confederacy had created a constitution and elected its first and only president, Jefferson Davis. The Constitution of the Confederate States of America was the supreme law of the Confederate States of America, as adopted on March 11, 1861 and in effect through the conclusion of the American Civil War. The Confederacy also operated under a Provisional Constitution from February 8, 1861 to March 11, 1861. In regard to most articles of the Constitution, the document is a word-for-word duplicate of the United States Constitution. The original, hand-written document is currently located in the University of Georgia archives at Athens, Georgia. The major differences between the two constitutions was the Confederacy's greater emphasis on the rights of individual member states, and an explicit support of slavery. Fort Sumter and the Beginning of the War. Several federal forts were seized and converted to Confederate strongholds. By the time of Lincoln's inauguration only two major forts had not been taken. On April 11, Confederate General P. G. T. Beauregard demanded that Union Major Robert Anderson surrender Fort Sumter in Charleston, South Carolina. Sumter had a strategic position on an island defending Charleston's harbor. The supplies of the besieged forts would only last a few weeks. The Union sent ships to resupply the fort, but they were held off by Confederate ships. Beauregard's troops surrounded the fort and opened fire. A tremendous cannon firefight remarkably claimed no casualties. By April 14, Anderson was forced to surrender the fort. The first casualties of the War occurred after the surrender: while the fort flag was being lowered, a Union cannon misfired. The next day, President Lincoln declared that the US faced a rebellion. Lincoln called up state militias and requested volunteers to enlist in the Army. In response to this call and to the surrender of Fort Sumter, four more states seceded; Virginia, Arkansas, Tennessee, and North Carolina. The Civil War had begun. Each side determined its strategies. The Confederate leadership felt that its army only needed to defend itself to gain independence. By its tactical strengths and its material shortages, it created what Jefferson Davis named an "offensive defensive" strategy. It would strengthen its defense posture, when conditions were right, by occasional offensive strikes into the North. However, three people who had important roles in Confederate plans had different strategies. While President Davis argued for a solely defensive war, General Robert E. Lee asserted they had to fight the Union head on, and General Thomas Jackson claimed they needed to invade the Union's important cities first and defeat the enemy to reclaim the cities. The strategy of aging Union General Winfield Scott became popularly known as the "Anaconda Plan". Named for the South American snake that strangles its victims to death, the plan aimed to defeat the Confederacy by surrounding it on all sides with a blockade of Southern ports and the swift capture of the Mississippi River. First Battle of Bull Run and the Early Stages of the War. Four slave states remained in the Union; Delaware, Maryland, Kentucky, and Missouri. The four border states were all important, and Lincoln did not want them to join the Confederacy. Missouri controlled parts of the Mississippi River, Kentucky controlled the Ohio river, and Delaware was close to the important Pennsylvania city of Philadelphia. Perhaps the most important border state was Maryland. It was close to the Confederate capital, Richmond, Virginia, and the Union capital of Washington was located between pro-Confederate sections of Maryland and seceded Virginia. Lincoln knew he had to be cautious if he did not want these states to join the Confederacy. But after the Battle of Fort Sumter, all of these states except for Maryland joined the South. Both sides had strengths and weaknesses. The North had a greater population, more factories, more supplies, and more money than the South. The South had more experienced military leadership, better-trained armies, and the advantage of fighting on familiar territory. Robert E. Lee is an example of the leadership the South relied upon. Before the Civil War, President Lincoln asked him to lead the Union army. Even though Lee was himself against slavery, he followed the people of his home state of Virginia into succession. Support for secession and the war was not unanimous in the Confederacy, and all of the southern states provided substantial numbers of troops for the Union armies. Moreover, the presence of slavery acted as a drain of southern manpower, as adult males who might otherwise join the army were required to police the slaves. On July 21, 1861, the armies of General Beauregard and Union General Irvin McDowell met at Manassas, Virginia in the Battle of Bull Run. Here the North originally had the upper hand, but Confederate General Thomas Jackson and his troops blocked Northern progress. Jackson's men began to retreat but Jackson stayed, standing "as a stone wall" (he was hereafter nicknamed "Stonewall Jackson"). As Confederate reinforcements arrived, McDowell's army retreated in confusion and was totally defeated. Before this, the North had nurtured a hope of quick victory over the Confederacy. The loss killed that hope. Though the Confederates achieved victory, General Beauregard did not chase stragglers of the defeated Union Army. Angered by this, Davis replaced him with General Robert E. Lee. Northern general McDowell's defeat by Confederates caused his replacement by George McClellan. Early Southern victories raised the complete defeat of the Union. The Confederacy appointed two representatives to the United Kingdom and France. Both traveled to Europe on a British ship, the "RMS Trent". A Union Captain, Charles Wilkes, seized the Trent and forced the Confederate representatives to board the Union ship. This seizure violated the neutrality of the United Kingdom. The British demanded apologies, and Lincoln eventually complied, even releasing the Confederate representatives. If he had failed to do so, the United Kingdom would have had an excuse to join with the Confederacy against the Union. Factories in the North of England depended upon Confederate cotton, and their neutrality was not assured. Technology. The Civil War was affected by technological innovations that changed the nature of battle. The most lethal change was the introduction of rifling to muskets. In previous wars, the maximum effective range of a musket was between 70 to 110 meters. Muskets, which were smooth bore firearms, weren't accurate beyond that. Tactics involved moving masses of troops to musket range, firing a volley, and then charging the opposing force with the bayonet, which is a sword blade attached to a firearm. However, a bullet from an aimed rifled musket could hit a soldier more than 1300 meters away. This drastically improved any defense. Massed attacks were easier to stop from a longer distance. The standardization of the rifle during the revolutionary war was extended to these new armaments, and to other military supplies. Some other key changes on land dealt with logistics -- the art of military supply -- and communications. By 1860, there were approximately 30,000 miles of railroad track, mostly in the Northern states. The railroads meant that supplies need not be obtained from local farms and cities, and that armies could operate for extended periods of time without fear of starvation. The advances in food preservation created during the Napoleonic Wars brought a wider variety of food to the soldier. In addition, armies could be moved across the country within days, without marching. Doctors could move to the wounded. The telegraph is the third of the key technologies that changed the nature of the war. Washington City and Richmond, the capitals of the two opposing sides, could stay in touch with commanders in the field, passing on updated intelligence and orders. President Lincoln used the telegraph frequently, as did his chief general, Halleck, and field commanders such as Grant. At sea, the greatest innovation was the introduction of ironclad warships. In 1862, the Confederate Navy built the CSS "Virginia" on the half-burned hull of the USS Merrimack. This ship, with iron armor, was impervious to cannon fire that would drive off or sink a wooden ship. The Virginia sank the U.S. frigate "Cumberland". It might have broken the blockade of the Federal fleet if it had not been for the arrival of the ironclad USS "Monitor", built by Swedish-American John Ericsson. The two met in May 1862 off Hampton Roads, Virginia. The battle was a draw, but this sufficed for the Union to continue its blockade of the Confederacy. The Virginia'had retreated into a bay where it could not be of much use, and the Confederacy later burned it to prevent Union capture. The U.S. Civil War introduced the first American railroad artillery; a successful submarine; a "snorkel" breathing device; the periscope for trench warfare; field trenches, land-mine fields, and wire entanglements, as battles began to take place for days at a time; American use of flame throwers and naval torpedoes; aerial reconnaissance, using hot-air balloons and cameras, and antiaircraft fire; resultant camouflage and blackouts; repeating rifles; telescopic sights for rifles for the aid of snipers, fixed ammunition, and long-range rifles for general use; electronic exploding bombs and torpedoes; revolving gun turrets on boats; and a workable machine gun. As part of the organization of men and materiel, the Civil War introduced foreign social innovations such as incorporation of female and civilian support in the Northern Sanitation Fairs, an organized medical and nursing corps with bandages, opium, and other anesthetics, hospital ships, and an army ambulance corps. To supply newspapers and magazines, with their sophisticated new engraving devices, there arose a wide-range corps of press correspondents in war zones. New aids in communication included the bugle call, "Taps," and other new calls, and the wigwag signal code in battle. To enable the federal prosecution of the war, the North inaugurated American conscription, legal voting for servicemen, The U.S. Secret Service, the income, withholding, and tobacco (cigarette) taxes, and the Medal of Honor. The Southern forces created a Confederate Department of Justice. The North created the first U.S. Navy admiral. Both sides commissioned Army Chaplains. The North commissioned African-American fighters, and its first African-American U.S. Army Officer, Major M.R. Delany. Shiloh and Ulysses Grant. While Union military efforts in the East were frustrated and even disastrous, the war west of the Appalachians developed differently, resulting in the first significant battlefield successes for the North. On the border between the Union and Confederacy, Kentucky was divided in its sentiments toward the two sides and attempted political neutrality. By the autumn of 1861, the Kentucky state government decided to support the Union, despite its being a slave state. Its indecision and the divided loyalties of its population directed the course of military operations in the West; neither North or South wished to alienate Kentucky. Below the confluence of the Ohio and Mississippi Rivers where the Kentucky, Tennessee and Missouri borders come together, Union Brigadier General Ulysses S. Grant, under command of Major General Henry W. Halleck, conducted a series of operations that would bring him national recognition. It was just across the Mississippi from Kentucky in Columbus, Missouri that Grant fought his first major battle. The western campaigns continued into 1862 under Halleck's overall direction with Grant continuing into Western Tennessee along the Mississippi. In February, Grant attacked and captured the Tennessean Fort Donelson, providing a significant victory for the North. About two months after the victory at Fort Donelson, Grant fought an even more important battle at Shiloh. Confederate generals A. S. Johnston and P. G. T. Beauregard made a surprise attack on the Union army. Though the initial attack was successful, the Union made a counter-attack and the Confederates were defeated. After the Union took Fort Donelson, Grant wanted to push onto into Charleston and Memphis, perhaps gaining control of the Eastern railroad and supply line. But General Helleck vetoed their proposal. Grant's troops killed Confederate General Albert Johnston and defeated the Confederate troops, but at a steep price. Approximately thirteen thousand Union soldiers and eleven thousand Confederate soldiers died, and Grant lost a chance of capturing the West quickly. Peninsular Campaign. General Stonewall Jackson was nearing Washington. To prevent Jackson from invading, Union General George McClellan left over fifty thousand men in Washington. Yet Jackson's threat was deceptive, as he did not even have five thousand men in his army. McClellan's unnecessary fear forced him to wait over half a year before continuing the war in Virginia, allowing enough time for the Confederates to strengthen their position and earning him the nickname "Tardy George. Jackson's deception had a further effect in the Peninsular Campaign, the Union attempt to take the Confederate capital Richmond "without" the aid of the force remaining in Washington. (The Union strategy for a quick end to the war was capturing Richmond, which was close to Washington.) In early April 1862, McClellan's troops began the Campaign, traveling over sea to the peninsula formed by the mouths of the York and James Rivers. This spit of land included Yorktown and Williamsburg and led straight to Richmond. By late May, McClellan was a few miles from Richmond, when Robert E. Lee took control of one of the Confederate Armies. After several victorious battles, it seemed as if McClellan could march to Richmond. But he refused to attack without reinforcements, which he saw as necessary to defeat Jackson's illusory troops. The forces he wanted were instead defending Washington. During the last week of June, Confederate General Robert E. Lee started the Seven Days' Battles that forced McClellan to retreat. By July, McClellan had lost over fifteen thousand men: there was little consolation in the fact that Lee had lost even more. Other important skirmishes occurred in the course of the Peninsular Campaign. Flag Officer David Farragut of the Union Navy easily took control of the Mississippi River when he captured the key port of New Orleans in April, providing a key advantage to the Union and depriving the Confederacy of the river. The North raised a blockade around the ports of the South, cutting off dry goods such as shoes and vastly increasing inflation. (Although the Confederates produced raw materials, they did not have the industrial wherewithal to finish them -- for example, the cotton mills in the North and abroad -- or the railroads to fully distribute them.) Second Bull Run and Antietam. A new Union Army was organized at the same time under General John Pope. Pope attempted to join his army with McClellan's to combine their strengths. Stonewall Jackson headed this off by surrounding Pope's Army in Manassas, which the North called the Second Battle of Bull Run. Both sides fought on August 29, and the Confederates won against a much larger Union force. Pope's battered Army did eventually combine with McClellan's. But the Second Battle of Bull Run had encouraged General Lee to invade Maryland. In Sharpsburg, Maryland, McClellan and Lee led their armies against each other. On September 17, 1862, the Battle of Antietam (named for a nearby creek) led to the deaths of over ten thousand soldiers from each side; no other one-day battle led to more deaths in one day. This day is called "Bloodiest day of American History". McClellan's scouts had found Lee's battle plans with a discarded packet of cigars, but he did not act on the intelligence immediately. The Union technically won the Pyrrhic victory; McClellan lost about one-sixth of his Army, but Lee lost around one-third of his. Even though they could march and end the war, McClellan didn't go forward because he thought he's already lost too many soldiers. This was the victory needed for Lincoln's Emancipation Proclamation, so that it did not appear as an act of desperation. The Emancipation Proclamation. General McClellan seemed too defensive to Lincoln, who replaced McClellan with General Ambrose Burnside. Burnside decided to go on the offensive against Lee. In December 1862, at Fredricksburg, Virginia, Burnside's Army of the Potomac assaulted built-up Confederate positions and suffered terrible casualties to Lee's Army of Northern Virginia. The Federal superiority in numbers was matched by Lee's use of terrain and modern firepower. "Burnside's Slaughter Pen" resulted in over ten thousand Union casualties, as the North used Napoleonic tactics against the South's carbines. Burnside then again attempted to capture Richmond, but was foiled by winter weather. The "Mud March" forced the Army of the Potomac to return to winter quarters. President Lincoln liked men who did not campaign on the abolition of slavery. He only intended to prevent slavery in all new states and territories. On the 22nd of August, 1862, Lincoln was coming to the decision that abolishing slavery might help the Union, in a letter from that time he wrote "My paramount object in this struggle is to save the Union, and is not either to save or destroy slavery. If I could save the Union without freeing any slave, I would do it; and if I could save it by freeing all the slaves, I would do it; and if I could do it by freeing some and leaving others alone, I would also do that.". Doing so would especially disrupt the Confederate economy. In September, 1862, after the Battle of Antietam, Lincoln and his Cabinet agreed to emancipate, or free, southern slaves. On January 1, 1863, Lincoln issued the Emancipation Proclamation, which declared all slaves in rebel states "forever free." The constitutional authority for the Emancipation Proclamation cannot be challenged. The Proclamation did not abolish slavery everywhere; it was restricted to states "still in rebellion" against the Union on the day it took effect. The Proclamation, technically, was part of a military strategy against states that had rebelled; this was to prevent internal conflict with the border states. Still, all the border states except Kentucky and Delaware had abolished slavery on their own. Naturally, the proclamation had no way of being enforced: the Executive in the form of military action was still trying to force the Confederacy to rejoin. Nonetheless, many slaves who had heard of the Proclamation escaped when Union forces approached. The Proclamation had another profound effect on the war: it changed the objective from forcing the Confederacy to rejoin the Union to eliminating slavery throughout the United States. The South had been trying too woo Great Britain (which relied on the South's agricultural exports, especially cotton, for manufacturing) into an alliance; now all hopes for one were eliminated. Great Britain was firmly against the institution of slavery, and it had been illegal throughout the British Empire since 1833. In fact, some slaves freed via the Underground Railroad were taken to Britain, since it was safe from bounty hunters. (Canada was too close to the U.S. for some). Although the Union did not at first accept black freedmen for combat, it hired them for other jobs. When troops became scarce, the Union began enlisting blacks. At the end of the war, the 180,000 enlisted blacks made up about 10% of the Union Army, and 29500 enlisted blacks to Navy. Until 1864, the South refused to recognize captured black soldiers as prisoners of war, and executed several of them at Fort Pillow as escaped slaves. Lincoln believed in the necessity of black soldiers: in August 1864, he said if the black soldiers of the Union army all joined the Confederacy, "we would be compelled to abandon the war in three weeks." Fredericksburg and Chancellorsville. In 1863, Lincoln again changed leadership, replacing Burnside with General Joseph Hooker. Hooker had a reputation for aggressiveness; his nickname was "Fighting Joe". From May 1 to May 4, 1863, near Chancellorsville, Virginia, General Lee, again outnumbered, used audacious tactics — he divided his smaller force in two in the face of superior numbers, sending Stonewall Jackson to the Union's flank, and defeated Hooker. Again, the Confederacy won, but at a great cost. Shortly after the battle of Chancellorsville, Stonewall Jackson was accidentally shot by Confederate soldiers who didn't recognize him in the poor evening light, dying soon after. Vicksburg. The North already held New Orleans. If it could control the entire Mississippi River, it could divide the Confederacy in two, making Confederate transportation of weapons and troops more difficult. Vicksburg and Fort Hudson were major Confederate ports. General Scott's "Anaconda Plan" was based on gaining control of the Mississippi. The city of Vicksburg, Mississippi, was located on high bluffs on the eastern bank of the river. At the time, the Mississippi River went through a 180-degree U shaped bend by the city. (It has since shifted course westward and the bend no longer exists.) Guns batteries there prevented Federal steamboats from crossing. Vicksburg was also on one of the major railroads running east-west through the Confederacy. Vicksburg was therefore a key point under Confederate control. Major General Ulysses Grant marched on land from Memphis, Tennessee, while General William Tecumseh Sherman and his troops traveled by water. Both intended to converge on Vicksburg. Both failed, at least for the time being. In December, 1862, Grant's supply line was disrupted, and Sherman had to attack alone. Since Vicksburg had not fallen to a frontal assault, the Union forces made several attempts to bypass Vicksburg by building canals to divert the Mississippi River, but these failed. Grant decided to attack Vicksburg again in April. Instead of approaching from the north, as had been done before, his army approached Vicksburg from the south. Grant's Army of Tennessee crossed from the West bank to the East at Big Bluff on April 18, 1863. Then, in a series of battles, including Raymond and Champion's Hill, defeated Southern forces coming to the relief of Confederate general Pemberton. Sherman and Grant together besieged Vicksburg. Two major assaults were repelled by the defenders of Vicksburg, including one in which a giant Union land mine was set off under the Confederate fortifications. From May to July, Vicksburg remained in Confederate hands, but on July 3, 1863, one day before Independence Day, General Pemberton finally capitulated. Thirty thousand Confederates were taken prisoner, but released after taking an oath to not participate in fighting the United States unless properly exchanged (a practice called parole). This victory cut the Confederate States in two, accomplishing one of the Union total war goals. Confederate forces would not be able to draw on the food and horses previously supplied by Texas. This victory was very important, giving the Union control of the whole Mississippi River and effectively splitting the Confederacy. Confederate forces were now deprived of food and supplies from Texas. Gettysburg. Background. At the same time as the opening of the Vicksburg Campaign, General Lee decided to march his troops into Pennsylvania. He had three reasons for doing this. He intended to win a major victory on Northern soil, increasing Southern morale, encouraging Northern peace activists sympathetic to the South (the "Copperheads"), and increasing the likelihood of political recognition by England and France. His hungry, poorly shod army could raid supplies from the North, reducing the burden on the Confederate economy. And he intended to encroach upon the Northern capital, forcing the recall of Federal troops from the Western Theater and easing some of the pressure on Vicksburg. Keeping the Blue Ridge Mountains between him and the Federal army, Lee advanced up the Shenandoah Valley into West Virginia and Maryland before finally marching into South-Central Pennsylvania. Meanwhile, the Union forces moved north on roads to Lee's east, without the latter's knowledge. His cavalry commander and chief scout, Jeb Stuart, had launched a raid eastward to "ride around" the Union army. On July first, 1863, Confederate Division Leader Henry Heth's soldiers ran into John Buford's Federal cavalry unit west of the city of Gettysburg. Buford's two brigades held their ground for several hours, until the arrival of the Union 1st Corps, and then withdrew through the town. The Confederates occupied Gettysburg, but by then the Union forces had formed a strong defensive line on the hills south of the town. The Battle. For the next three days, the Confederate Army of Northern Virginia faced the Union Army of the Potomac, now under the command of General George G. Meade, a Pennsylvanian who replaced Hooker, who had resigned as commander. (Hooker was given a corps command in the Army of the Cumberland, then in eastern Tennessee, where he performed satisfactorily for the remainder of the war.) South of Gettysburg are high hills shaped like an inverted letter "J". At the end of the first day, the Union held this important high ground, partially because the Confederate left wing had dawdled moving into position. One July 2, Lee planned to attack up Emmitsburg Road from the south and west, hoping to force the Union troops to abandon the important hills and ridges. The attack went awry, and some Confederate forces, including Law's Alabama Brigade, attempted to force a gap in the Federal line between the two Round Tops, dominant heights at the extreme southern end of the Union's fish hook-shaped defensive line. Colonel Joshua Lawrence Chamberlain, commander of the 20th Maine Regiment, anchored this gap. He and the rest of his brigade, commanded by Colonel Strong Vincent, held the hill despite several hard-pressed attacks, including launching a bayonet charge when the regiment was low on ammunition. Meanwhile, north of the Round Tops, a small ridge immediately to the west of the Federal line drew the attention of Union General Daniel Sickles, a former New York congressman, who commanded the Third Corps. He ordered his corps to advance to the peach-orchard crested ridge, which led to hard fighting around the "Devil's Den," Wheatfield, and Peach Orchard. Sickles lost a leg in the fight. Pickett's Charge. On the third day of the Battle of Gettysburg, Lee decided to try a direct attack on the Union and "virtually destroy their army." Putting Lieutenant General James Longstreet in charge of the three-division main assault, he wanted his men, including the division of Major General George Pickett, to march across a mile and a half up a gradual slope to the center of the Union line. Lee promised artillery support, but any trained soldier who looked across those fields knew that they would be an open target for the Union soldiers--much the reverse of the situation six months before in Fredericksburg. However, the choice was either to attack or withdraw, and Lee was a naturally aggressive soldier. By the end of the attack, half of Longstreet's force was dead, wounded or captured and the position was not taken. George Pickett never forgave Lee for "slaughtering" his men. Pickett's Charge, called the "High Water Mark of the Confederacy," was practically the last hope of the Southern cause at Gettysburg. Aftermath &amp; The Gettysburg Address. Lee withdrew across the Potomac River. Meade did not pursue quickly, and Lee was able to reestablish himself in Virginia. He offered to Confederate President Jefferson Davis to resign as commander of the Army of Northern Virginia, saying, ""Everything, therefore, points to the advantages to be derived from a new commander, and I the more anxiously urge the matter upon Your Excellency from my belief that a younger and abler man than myself can readily be attained." Davis did not relieve Lee; neither did Lincoln relieve Meade, though he wrote a letter of censure, saying "Again, my dear general, I do not believe you appreciate the magnitude of the misfortune involved in Lee's escape. He was within your easy grasp, and to have closed upon him would, in connection with our other late successes, have ended the war. As it is, the war will be prolonged indefinitely." The battle of Gettysburg lasted three days. Both sides lost nearly twenty-five thousand men each. After Gettysburg, the South remained on the defensive."' On November 19, 1863 Lincoln delivered his most famous speech in the wake of this battle. The Gettysburg Address is often cited for its brevity (it followed a two-hour speech by Edward Everett) and its masterful rhetoric. As with other early Republican documents, it placed its justification in the Founding Fathers. Unlike them, it did not place the justification of emancipation in the Constitution, but in the Declaration of Independence: "All men are created equal." Black Americans and the Civil War. The view of the Union towards blacks had changed during the previous two years. At the beginning of hostilities, the war was seen as an effort to save the Union, not free slaves. Several black slaves who reached Federal lines were returned to their owners. This stopped when Major General Benjamin F. Butler, a New Jersey lawyer and prominent member of the Democratic party, announced that slaves, being the property of persons in rebellion against the United States, would be seized as "contraband of war" and the Fugitive Slave Act could not apply. "Contrabands" were, if not always welcome by white soldiers, not turned away. However, as the struggle grew more intense, abolition became a more popular option. Frederick Douglas, a former slave, urged that the war aim of the Union include the emancipation of slaves and the enlistment of black soldiers in the Union Army. This was done on a nationwide basis in 1863, though the state of Massachusetts had raised two regiments (the 54th and 55th Massachusetts) before this. The 54th Massachusetts Regiment was the first black regiment recruited in the North. Col. Robert Gould Shaw, the 25 year old son of very wealthy abolitionist parents, was chosen to command. On May 28, the well equipped and drilled 54th paraded through the streets of Boston and then boarded ships bound for the coast of South Carolina. Their first conflict with Confederate soldiers came on July 16, when the regiment repelled an attack on James Island. But on July 18 came the supreme test of the courage and valor of the black soldiers; they were chosen to lead the assault on Battery Wagner, a Confederate fort on Morris Island at Charleston. In addressing his soldiers before leading them in charge across the beach, Colonel Shaw said, "I want you to prove yourselves. The eyes of thousands will look on what you do tonight." While some blacks choose to join the military fight others fought by other means. An American teacher named Mary S. Peake worked to educate the freedmen and "contraband". She spent her days under a large oak tree teaching others near Fort Monroe in Virginia. (This giant tree is now over 140 years old and called Emancipation Oak). Since Fort Monroe remained under Union control this area was some what of a safe location for refugees and runaways to come to. Soon Mary began teaching in the Brown Cottage. This endeavor, sponsored by the American Missionary Association, became the basis from which Hampton University would spawn. Mary's school would house around 50 children during the day and 20 adults at night. This remarkable American died from tuberculosis on Washington's birthday in 1862. Confederate President Jefferson Davis reacted to the raising of black regiments by passing General Order No. 111, which stated that captured black Federal soldiers would be returned into slavery (whether born free or not) and that white officers who led black soldiers would be tried for abetting servile rebellion. The Confederate Congress codified this into law on May 1, 1863. President Lincoln's order of July 30, 1863 responded: Eventually the Federal forces had several divisions' worth of black soldiers. Their treatment was not equal to white soldiers: at first, for example, black privates were paid $10 a month, the same as laborers, while white privates earned $13 a month. In addition, blacks could not be commissioned officers. The pay difference was settled retroactively in 1864. The issue of black prisoners of war was a continual contention between the two sides. In the early stages of the war, prisoners of war would be exchanged rank for rank. However, the Confederates refused to exchange any black prisoner. The Union response was to stop exchanging any prisoner of war. The Confederate position changed to allowing blacks who were born free to be exchanged, and finally to exchange all soldiers, regardless of race. By then, the Federal leadership understood that the scarcity of white Confederates capable of serving as soldiers was an advantage, and there were no mass exchanges of prisoners, black or white, until the Confederate collapse. Chickamauga and Chattanooga. In September 1863, Union Major General William Rosecrans decided to attempt the takeover of Chattanooga, a Confederate rail center in the eastern part of Tennessee. Controlling Chattanooga would provide a base to attack Georgia. The Confederates originally gave up Chattanooga, thinking that they could launch a devastating attack as the Union Army attempted to take control of it. Rosecrans did not, in the end, fall into such a trap. However, on November 23, 1863, the Union and Confederate Armies met at Chickamauga Creek, south of Chattanooga, upon which a rail line passed into Georgia. The battle of Chickamauga was a Confederate victory. The Army of the Cumberland was forced to withdraw to Chattanooga, but Union General George Thomas, "the Rock of Chickamauga," and his troops prevented total defeat by standing their ground. After Rosecrans withdrew to Chattanooga, the Confederates under General Braxton Bragg decided to besiege the city. Rosecrans was relieved of command; Lincoln's comment was that he appeared "stunned and confused, like a duck hit on the head." Meanwhile, by great effort, the Federal forces kept a "cracker line" open to supply Chattanooga with food and forage. Ulysses Grant replaced Rosecrans. Grant's forces began to attack on November 23, 1863. On November 24 came the Battle of Lookout Mountain, an improbable victory in which Union soldiers, without the initiative of higher command, advanced up this mountain, which overlooks Chattanooga, and captured it. One of the authors of this text had an ancestor in the Confederate forces there; his comment was when the battle started, he was on top of the hill throwing rocks at the Yankees, and when it was over, the Yankees were throwing rocks at him. By the end of November, Grant and his troops had pushed the Confederates out of East Tennessee and begun operations in Georgia. Ulysses Grant As General-in-Chief. Lincoln recognized the great victories won by Ulysses Grant. In March, 1864, the President made Grant the general-in-chief of Union Forces, with the rank of Lieutenant General (a rank only previously held by George Washington). Grant decided on a campaign of continual pressure on all fronts, which would prevent Confederate forces from reinforcing each other. He went east and made his headquarters with General Meade's Army of the Potomac (although Grant never took direct command of this Army). The Army of the Potomac's chief mission would be to whittle down the manpower of the Army of Northern Virginia, Lee's army. In May 1864, the two sides met in Virginia near site of the previous year's Battle of Chancellorsville. The terrain was heavily wooded and movement to attack or reinforce was particularly difficult. During the Battle of the Wilderness, the Union lost eighteen thousand soldiers, while the Confederates lost eleven thousand. Nevertheless, the Union pushed on. The two Armies fought each other again at Spotsylvania Court House and at Cold Harbor. In each case, the Union again lost large numbers of soldiers. Grant then hatched a plan to go "around" rather than through the Confederate Army in order to capture Richmond. At the last second, due to a hesitation by Major General "Baldy" Smith, the Army of Northern Virginia blocked the Union troops at Petersburg. Grant then decided to siege the city (and Lee's forces) and force it to surrender; if Lee could not move, he could not help other Confederate armies. The siege took almost one year. The Georgia Campaign and Total War. Battles for Atlanta. This victory had a significant effect on the election of 1864. Without it, there might have been more support for his Copperhead opponent General McClellan. The March to the Sea. The ultimate Union strategy emerged with six parts: blockade the Confederate coastlines, preventing trade; free the slaves, destroying the domestic economy; disconnect the Upper South from the Deep South by controlling the Mississippi River; further split the Confederacy by attacking the Southeast coast (Georgia, South Carolina, and North Carolina), denying access to foreign supply; capture the capital of Richmond, which would severely incapacitate the Confederacy; and engage the enemy everywhere, weakening the army through attrition. If Richmond had indeed been captured quickly and the war had ended within a few months, the Plantation system and slavery would probably not have changed significantly. Because the South was fighting predominately in its own territory, primarily rural farmland, its soldiers could take or force food and support from the people around them. After the unsuccessful Union attacks in Virginia, Lincoln began to think about the Emancipation Proclamation, and the Union changed its strategy from a quick capture of Richmond to the destruction of the South through "total war". In total war, an invading army destroys both military and non-combatant resources important to war. It can involve attacks on civilians or the destruction of civilian property. General William Sherman used total war in his March to the Sea in November and December in 1864. Once Atlanta was taken, General Sherman and four army corps disconnected themselves from any railroad or telegraphic communications with the Union and headed through the state of Georgia. Their objective was Savannah, Georgia, a major seaport. Sherman's strategy was to inflict as much damage on the civilian population of Georgia as possible, short of killing people. To accomplish this, he issued orders to "forage liberally on the country." Many of his soldiers saw this as a license to loot any food or valuable property they could. Sherman officially disapproved of this. Sherman's army carved a path of destruction 300 miles long and over 60 miles wide from Atlanta to the coastal city of Savannah. It destroyed public buildings and railroad tracks wherever it went. Troops heated railroad rails to white heat and twisted them around the trees, creating "Sherman's neckties." Sherman's strategy separated his forces from the main body of the Union army, yet maintained the men with food and weapons. It not only aided his regiments without supply lines -- Southern destruction of supply lines had previously halted Northern advances -- but destroyed supply caches for Confederate forces in the area as well. But this destruction combined with Southern army raids to throw non-combatants into starvation. On his way to Savannah, Sherman did not burn down every town he passed through, choosing to spare some such as Madison, Georgia for political reasons. The Confederate forces were unable to take on Sherman's forces, and evacuated, leaving behind large amounts of supplies in the city of Savannah. Undefended, the historic city of Savannah surrendered to Sherman, and it was spared. He reached the city of Savannah on December 24, 1864, and telegraphed President Lincoln "I present to you the city of Savannah as a Christmas present." Moving through the Carolinas. Sherman's forces then moved north into South Carolina, while faking an approach on Augusta, Georgia; the general's eventual goal was to coordinate his forces with those of General Grant in Virginia and entrap and destroy Lee's Army of Northern Virginia. The pattern of destruction by the Union soldiers continued, often with a more personal feeling of vengeance. A Federal soldier said to his comrades, "Here is where treason began and, by God, here is where it will end!" On February 17, 1865, Sherman's forces reached Columbia, the capital of South Carolina. After a brief bombardment, the city surrendered. However, a large stock of whiskey was left behind as the Confederates retreated. Drunken soldiers broke discipline; convicts were let loose from the city jail, and somehow fires broke out, destroying much of the city. Hood's Invasion of Tennessee and the Battle of Nashville. Spring Hill. The battle of Spring Hill was fought on November 29, 1864, at Spring Hill, Tennessee. The Confederates attacked the Union as it retreated from Columbia. The Confederates were not able to inflict significant damage to the retreating Union force. So the Union Army was still able to make it safely north to Franklin during the night. The following day the Confederates decided to follow the Union and attack a much more fortified group at the Battle of Franklin. This did not prove to be a wise decision, as the Confederates suffered many casualties. Franklin. The Battle of Franklin was fought on November 30, 1864 at Franklin, Tennessee. This battle was a devastating loss for the Confederate Army. It detrimentally shut down their leadership. Fourteen Confederate Generals were extinguished with 6 killed, 7 wounded and 1 captured. 55 Regimental Commanders were casualties as well. After this battle the Confederate Army in this area was effectively handicapped. Nashville. In one of the decisive battles of the war, two brigades of black troops helped crush one of the Confederacy's finest armies at the Battle of Nashville on December 15-16, 1864. Black troops opened the battle on the first day and successfully engaged the right of the rebel line. On the second day Col. Charles R. Thompson's black brigade made a brilliant charge up Overton Hill. The 13th US Colored Troops sustained more casualties than any other regiment involved in the battle. Fort Pillow. The Battle of Fort Pillow was fought on was fought on April 12, 1864, at Fort Pillow on the Mississippi River at Henning, Tennessee. The battle ended with a massacre of surrendered Union African-American troops under the direction of Confederate Brigadier General Nathan Bedford Forrest. The End of the Confederacy. The Siege of Petersburg. The Siege of Petersburg, also known as The Richmond Petersburg Campaign, began on June 15, 1864 with the intent by the Union Army to take control of Petersburg which was Virginia's second largest city and the supply center for the Confederate capital at Richmond. The campaign lasted 292 days and concluded with the occupation of Union forces on April 3, 1865. Thirty-two black infantry and cavalry regiments took part in the siege. First Battle of Deep Bottom. The First Battle of Deep Bottom is also known as Darbytown, Strawberry Plains, New Market Road, and Gravel Hill. It was part of The Siege of Petersburg, and was fought July 27-29, 1864, at Deep Bottom in Henrico County, Virginia. The Crater. The Battle of the Crater was part of the Siege of Petersburg and took place on July 30, 1864. The battle took place between the Confederate Army of Northern Virginia and the Union Army of Potomac. The battle was an unusual attempt by the Union to penetrate the Confederate defenses south of Petersburg, VA. The battle showed to be a Union disaster. The Union Army went into battle with 16,500 troops, under the direct command of Ulysses S. Grant; the Confederate Army was commanded by Robert E. Lee and entered battle with 9,500 troops. Pennsylvania miners in the Union general Ambrose E. Burnside's Ninth Corps, worked for several weeks digging a long tunnel, and packing it with explosives. The explosives were then detonated at 3:15 on the morning of July 30, 1864. Burnside originally wanted to send a fresh division of black troops against the breach, but his superiors, Ulysses S. Grant, ruled against it. The job, chosen by short straw, went to James H. Ledlie. Ledlie watched from behind the lines as his white soldiers, rather than go around, pile into the deep crater, which was 170 feet long, 60 feet across, and 30 feet deep. They were not able to escape making the Union soldiers easy targets for the Confederates. The battle was marked by the cruel treatment of black soldiers who took part in the fight, most of them were captured and murdered. The battle ended with a confederate victory. The Confederacy took out 3,798 Union soldiers, while the Union were only able to defeat 1,491 Confederate soldiers. The United States Colored Troops suffered the most with their casualties being 1,327 which would include 450 men being captured. Second Deep Bottom. The Second Battle of Deep Bottom was fought August 14-20, 1864, at Deep Bottom in Henrico County, Virginia; it was part of the Siege of Petersburg. The battle is also known as Fussell's Mill, Kingsland Creek, White's Tavern, Bailey's Creeks, and Charles City Road. General Winfield Scott Hancock came across the James River at Deep Bottom where he would threaten Richmond, Virginia. This would also cause the Confederates to leave Peterburgs, Virginia and the trenches and Shenandoah Valley. Appomattox. Sherman did not stop in Georgia. As he marched North, he burnt several towns in South Carolina, including Columbia, the capital. (Sherman's troops felt more anger towards South Carolina, the first state to secede and in their eyes responsible for the war.) In March 1865, Lincoln, Sherman, and Grant all met outside Petersburg. Lincoln called for a quick end to the Civil War. Union General Sheridan said to Lincoln, "If the thing be pressed I think Lee will surrender." Lincoln responded, "Let the thing be pressed." On April 2, 1865, the Confederate lines of Petersburg, Richmond's defense, which had been extended steadily to the west for 9 months, broke. General Lee informed President Davis he could no longer hold the lines; the Confederate government then evacuated Richmond. Lee pulled his forces out of the lines and moved west; Federal forces chased Lee's forces, annihilated a Confederate rear guard defense, and finally trapped the Army of Northern Virginia. General Lee requested terms. The two senior Army officers met each other near Appomattox Courthouse in Virginia on April 9th,1865. The men met at the home of Wilmer McLean. The gathering lasted about two and half hours. Grant offered extremely generous terms, requiring only that Lee's troops surrender and swear not to bear arms till the end of the War. This meeting helped to nearly end the bloodiest war in American history. General Sherman met with Confederate General Robert E. Lee to discuss the surrender of Confederate troops in the South. Sherman initially allowed even more generous terms than Grant. However, the Secretary of War refused to accept the terms because of the assassination of Abraham Lincoln by the Confederate John Wilkes Booth. By killing Lincoln at Ford's Theater, Booth made things worse for the Confederacy. Sherman was forced to offer harsher terms of surrender than he originally proposed, and General Johnston surrendered on April 26 under the Appomattox terms. All Confederate armies had surrendered by the end of May, ending the Civil War. Notable Raids. The Great Locomotive Chase resulted in the first Medal of Honor being issued. Morgan's Raid was a Confederate raid that went deep into Union territory. Besides the Fighting. Not all the important events of the Civil War took place on the battlefield. Petroleum Nasby. Operating under the pseudonym "Petroleum V. Nasby", journalist David Ross Locke gained a large amount of popularity by Union residents during the war, including by President Abraham Lincoln. "Petroleum V Nasby" was a mockery of Pro South Democrats, with his published letters being filled with misspellings, drunkenness, vitriol, bigotry, and a general desire to slack and grift his way to a comfy position as a postmaster. Domestic Affairs. On April 22, 1864, the U.S. Congress passed the Coinage Act of 1864 which mandates that the inscription "In God We Trust" be placed on all coins minted as United States currency. Dr. Rebecca Lee Crumpler becomes the first black woman to receive a medical degree. As far back as the 1850s Whig interests had introduced three bills to Congress: a homestead act, a Pacific railroad act, and grants to establish agricultural and technical colleges. These measures were seen as remedies for the depression of 1857. Southern interests had vetoed all of them. Now Republicans took advantage of a legislature free of slave interests. On May 20, 1862, the United States Congress passed the Homestead Act. Now any adult American citizen, or a person intending to become an American citizen, who was the head of a household, could qualify for a grant of 160 acres (67 hectares) of land by paying a small fee and living on the land continuously for five years. If a person was willing to pay $1.25 an acre, the time of occupation dwindled to six months. The Pacific Railway Acts of 1862 and 1864 enabled the United States Government to make a direct grant of land to railway companies for a transcontinental railroad, as well as a payment of $48,000 for every mile of track completed and lower-than-prime rate loans for any railway company who would build such a railway. The Central Pacific and the Union Pacific began to construct lines. The two railways finally met four years after the war, in Promontory Point, Utah, in 1869. The third major bill of these three, which established a land-grant university, is discussed below. The Draft. The federal government started a draft lottery in July 1863. Men could avoid the draft by paying $300, or hiring another man to take their place. This caused resentment among the lower classes as they could not afford to dodge the draft. On Monday, July 13, 1863, between 6 and 7 A.M., the Civil War Draft Riots began in New York City. Rioters attacked the draft offices, the Bull's Head Hotel on 44th Street, and more upscale residences near 5th Avenue. They lynched black men, burned down the Colored Orphan Asylum on 5th Avenue between 43rd and 44th Streets, and forced hundreds of blacks out of the city. Members of the 7th New York Infantry and 71st New York Infantry subdued the riot. Military Intelligence. Both the Union and the Confederacy operated intelligence gathering efforts during the Civil War. A number of Women conducted Espionage during the war. Harriot Tubman was one such spy for the Union. The Confederate Secret Service and the Confederate Signal Corps both conducted espionage for the Confederacy. The Union intercepted a number of Confederate cipher messages during the war. Indigenous People. While Lincoln proved to be instrumental in the emancipation of blacks, the Native Americans were not so lucky. Lincoln was responsible for the largest mass hanging in United States history. Thirty-eight Native Americans from the Santee Sioux tribe were hung on December 26, 1862. The US government failed to honor its treaties with the Indian Nations. They were supposed to supply the Indians with money and food for signing a treaty to turn over more than one million acres of land. Instead the agents kept the money and sold the food that was supposed to go the Indians to the white settlers. The food that was given to the Indians was spoiled and unfit for human consumption. Subsequently, the Indians went off the reservation in hunting parties to try to find suitable food. One of the Indian hunting groups took some eggs from a white settler's land and that caused this extreme government action. Authorities in Minnesota asked President Lincoln to order the execution of all 303 Indian males. However, Lincoln was afraid of how Europe would react so he tried to compromise. They would only execute those who were in the group. Lincoln also agreed to kill or remove every Indian from the state and provide Minnesota with 2 million dollars in federal funds. Ironically, he owed the Sioux only 1.4 million dollars for the land. Education. Land Grant Universities. In the Morrill Act of 1862, the government granted land to Union states to sell for funding educational institutions. This excluded the states which had seceded from the Union. The schools would teach military tactics, agriculture, and engineering. This answered the Republican campaign promise of 1860. These "Land Grant Universities" were proposed to spread small farm prosperity, as opposed to the large, inherited plantations, and to increase industrial innovations across a wider area. 1860's schoolhouses. In the 1860s, most schools were small, multiple grades were taught in one classroom at one time. Paper was expensive, and the more prosperous schools had students write their problems on individual student slates. Memorization was a common means of learning, and student knowledge was measured by oral recitation. Teachers often punished "bad children" with the dunce cap, a rap on a palm with a ruler, hitting or spanking, or even striking a child with a rod or a whip. Corporal punishment was seen as simply one way of enforcing obedience. Teacher and parent both generally agreed that obedience was the trait of good children. Literacy. Farming was still a major form of employment in America. It had been so since the first "semester", the time when students were allowed to be in school because the crops had been sown. Students worked in the fields during harvest time, and most left school for good to work on a farm. Abraham Lincoln himself, as a youth on the frontier, had had little schooling. Yet despite these brief periods of education, the reading levels were actually quite high. By the fifth grade students were sometimes reading books that we would consider college level, and Latin was still a part of many curricula. Academies. Academies during this time provided education for children between the ages of thirteen and twenty. These academies offered an array of classes. Most of the academies kept the boys and girls separate. There were also "seminaries", or private schools, which might cater to boys or girls. Girl's schools varied widely. Emily Dickinson's school, Amherst Academy, taught Mental Philosophy, Geology, Latin, and Botany. Some schools left girls idle, with not even what we would call physical education. Others taught non-intellectual, "feminine" skills such as deportment, needle craft, and perhaps arts and crafts. The Home Economics movement, inaugurated by Catherine Beecher, advocated teaching homemaking skills in school. It also promoted female physical education. In contrast, feminists such as Susan B. Anthony and Emma Willard, and reformers such as Jane Addams and Mary McLeod Bethune, wanted to expand women's education into the plane of men. These women helped establish the higher education institutions where women were able to take classes not otherwise offered to them. The first co-educational college was Oberlin College, established in 1833. The first all-women's college was Vassar College in 1861. Questions For Review. 1. What are the four principal causes of the Civil War? 2. Why did Sherman feel compelled to adopt the total war strategy in his March to the Sea? What are the advantages and disadvantages of this strategy? 3. The Morrill Act of 1862, the Homestead Act of 1862, and the Pacific Railroad Acts of 1862 and 1864: why did slave-holding Southern interests oppose their predecessors? What effect did they have? 

Start of Reconstruction. Congress passed the first Reconstruction Act on 2nd March, 1867. The South was now divided into five military districts, each under a major general. New elections were to be held in each state with freed male slaves being allowed to vote. The act also included an amendment that offered readmission to the Southern states after they had ratified the Fourteenth Amendment and guaranteed adult male suffrage. President Andrew Johnson immediately vetoed the bill but Congress re-passed the bill the same day. Andrew Johnson consulted General Ulysses S. Grant before selecting the generals to administer the military districts. Eventually he appointed John Schofield (Virginia), Daniel Sickles (the Carolinas), John Pope (Georgia, Alabama and Florida), Edward Ord (Arkansas and Mississippi) and Philip Sheridan (Louisiana and Texas). During the American Civil War, in which the nation decided how to handle the return of the seceded states and the status of the Freedmen (the newly freed slaves). Most scholars have accepted 1865-1877 as the boundaries for Reconstruction. The era itself was controversial and pitted various segments of American society against one another. Differing conceptions on how to restore the former Confederate States into the Union collided with diverse opinions concerning the status of African-Americans. The meaning of freedom itself was at stake in this crucial time period. The nascent Republican Party was divided between the mainstream which wanted a modicum of protection for blacks, and the Radicals, who wanted a thorough reorganization of Southern society. Conservative elements of this time period (in particular the Democrats) believed that the old order that governed relations between the states and between blacks and whites should remain intact. The bulk of African-Americans desired equal civil and political rights, protection of their person, and in many cases a redistribution of land and the break-up of the plantation system. These diverse perspectives enabled the period from 1865 to 1877 to be, in many ways, a grand experiment in interracial democracy, but the period was also dominated by tense political relations and a preponderant violence across the South. Definition. Reconstruction, in United States history, refers both to the period after the Civil War when the states of the breakaway Confederate States of America were reintegrated into the United States of America, and to the process by which this was accomplished. For victory in the American Civil War to be achieved, Northern moderate Republicans and Radical Republicans concurred that the Confederacy and its system of slavery had to be destroyed, and the possibility of either being revived had to be eliminated. Controversy focused on how to achieve those goals, and who would decide when they were achieved. The Radical Republicans held that reaching those goals was essential to the destruction of the Slave Power, and necessary to guaranteeing perpetual unity of the states, as well as a solution to the many problems of Freedmen. United States Senator Charles Sumner of Massachusetts, a Radical Republican, held that Congress should abolish slavery along with the Confederacy, extend civil and political rights to blacks, and educate black and white students together. The "moderates" claimed early success in achieving the goals by assurances that the former Confederates had renounced secession and abolished slavery. Most moderates, like Abraham Lincoln and Andrew Johnson, wanted suffrage for black army veterans but not other African Americans. Southern political leaders renounced secession and gave up slavery, but were angered in 1867 when their state governments were ousted by federal military forces, and replaced by Radical Republican governments made up of Freedmen, Carpetbaggers and Scalawags. Their primary instrument was the Black Codes (1865). These restricted the rights of Blacks and limited economic and educational opportunities. For example, there was very little, if any, employment available in the south. The Yankees may have won the war to end slavery, however the reconstruction did not benefit the African Americans who searched for employment. The Problem of Reconstruction. Reconstruction was the effort of rebuilding the South based on free labor instead of slave labor. The issue to Northern politicians was how it would be done. At the end of the Civil War, Congress proposed the Thirteenth Amendment, which sought to prohibit slavery. A state was not to gain re-admittance into the Union until it ratified the Amendment, but some states such as Mississippi were admitted despite failing to ratify. The Amendment became a part of the Constitution in December 6,1865. During this time many Northerners moved to the South to start new lives. Sometimes carrying their belongings in briefcases made of carpet, they were known by Confederate Southerners as "carpetbaggers." Confederate Southerners also had a derogatory name for southern whites who sided with the Republicans. They called them scalawags. The period just after the war also saw the rise of black codes, which restricted the basic human rights of freed slaves. Some of the more common codes seen were: race was dependent on blood, which meant if you had any amount of black blood in your body, you were considered black, freedmen could not get together unless accompanied by a white person, public restrooms and other facilities were segregated. This time in history was really volatile. Many racially motivated riots broke out all over the country. The hostilities the south held toward the north and the African Americans grew stronger and stronger. Ku Klux Klan. Ku Klux Klan (KKK) is the name of several past and present organizations in the United States that have advocated white supremacy, anti-Semitism, anti-Catholicism, racism, homophobia, anti-Communism and nativism. These organizations have often used terrorism, violence, and acts of intimidation, such as cross burning and lynching, to oppress African Americans and other social or ethnic groups. The first branch of the Ku Klux Klan was established in Pulaski, Tennessee, in May, 1866. A year later, a general organization of local Klans was established in Nashville in April, 1867. Most of the leaders were former members of the Confederate Army and the first Grand Wizard was Nathan Forrest, an outstanding general during the American Civil War. During the next two years Klansmen wearing masks, white cardboard hats and draped in white sheets, tortured and killed black Americans and sympathetic whites. Immigrants, who they blamed for the election of Radical Republicans, were also targets of their hatred. Between 1868 and 1870, the Ku Klux Klan played an important role in restoring white rule in North Carolina, Tennessee, and Georgia. The Klan's first incarnation was in 1866. Founded by veterans of the Confederate Army, its main purpose was to resist Reconstruction. It focused as much on intimidating "carpetbaggers" and "scalawags" as on putting down the freed slaves. The KKK quickly adopted violent methods. A rapid reaction set in. The Klan's leadership disowned violence as Southern elites saw the Klan as an excuse for federal troops to continue their activities in the South. The organization declined from 1868 to 1870 and was destroyed in the early 1870s by President Ulysses S. Grant's vigorous action under the Civil Rights Act of 1871 (also known as the Ku Klux Klan Act). At the end of the American Civil War, members of Congress attempted to destroy the white power structure of the Rebel states. The Freeman’s Bureau was established by Congress on March 3rd, 1865. The bureau was designed to protect the interests of former slaves. This included helping them to find new employment and to improve educational and health facilities. In the year that followed, the bureau spent $17,000,000 establishing 4,000 schools, 100 hospitals and providing homes and food for former slaves. Violence against African Americans started on the first days of Reconstruction and became more organized significant after 1867. Members of The Klan looked to frustrate Reconstruction. They also, tried to keep freedom in subjection. Terrorism dominated some counties and regions so, nighttime harassment, whippings, beatings, rapes, and murders became more common. The Klan's main purpose was political, even though, they tormented blacks who stood up for their rights. Active Republicans were the target of lawless night riders. When freedmen that worked for a South Carolina scalawag started to vote, terrorists went to the plantation and, in the words of a victim, "whipped every ... [black] man they could lay their hands on." Lincoln and Reconstruction. Lincoln firmly believed that the southern states had never actually seceded, because, constitutionally, they cannot. He hoped that the 11 states that seceded could be "readmitted" by meeting some tests of political loyalty. Lincoln began thinking about re-admittance early on. In his Proclamation of Amnesty and Reconstruction, which was issued in 1863, Lincoln established a simple process, hoping that Unionists would rise to political power rather than secessionalists. This plan would have granted presidential pardons to all southerners (save the political leaders at the time) who took an oath of future allegiance to the Union. Under Lincoln's plan, a state could be established as legitimate as soon as 10 percent of the voting population in the 1860 general election took this oath and a government was set up accepting the emancipation of the slaves. Rejecting Lincoln's Presidential reconstruction plan, radical Republicans in congress arguing that it was too lenient, passed the Wade-Davis bill in 1864, which proposed far more demanding terms. It required 50 percent of the voters to take the loyalty oath and allowed only those who could swear that they had never supported the confederacy to run for office or hold federal employment. Lincoln rejected this plan and pocket-vetoed the bill. In March 1865, Congress created a new agency, the Freedman's Bureau. This agency provided food, shelter, medical aid, help to find employment, education, and other needs for blacks and poor whites. The Freedman's Bureau was the largest scale federal aid relief plan at this time. It was the first large scale governmental welfare program. In 1864, his Vice Presidential running mate was the only Southern Senator to remain loyal to the Union - Andrew Johnson from Tennessee. After Lincoln was assassinated on April 14, 1865, and Johnson became President, the latter proved to be an obstacle to the Radical Republicans in Congress, who attempted to completely overhaul the Southern government and economy, which would have caused further tensions. In May, 1865, Johnson made his own proclamation, one that was very similar to Lincoln's. Offering amnesty to almost all Confederates who took an oath of allegiance to the Union, Johnson also reversed General Sherman's decision to set aside land for the express use of freed slaves. Not long after Johnson took office, all of the ex-Confederate states were able to be readmitted under President Johnson's plan. In 1866, Johnson vetoed two important bills, one that bolstered the protection that the Freedmen's Bureau gave to blacks and a civil rights bill that gave full citizenship to blacks. After realizing that if all of the Republicans, moderate and radical alike, united, they could overcome Johnson's vetoes, they soon passed the Civil Rights Act of 1866 and the Fourteenth Amendment. This amendment declared citizenship for all persons born in the United States and required the states to respect the rights of all US citizens. The Civil Rights Act outlawed the black codes that had been prevalent throughout the South. Over Johnson's vetoes, Congress passed three Reconstruction acts in 1867. They divided the southern states into five military districts under the control of the Union army. The military commander in charge of each district was to ensure that the state fulfilled the requirements of Reconstruction by ratifying the Fourteenth Amendment and by providing voting rights without a race qualification. Tennessee was not included in the districts because it had ratified the Fourteenth Amendment in 1866 and was quickly readmitted to the Union. In 1868, the House of Representatives impeached Andrew Johnson. Earlier, Congress had passed the Tenure of Office Act (over Johnson's veto), which required the President to dismiss officers only with the advice and consent of the Senate if he appointed them with the same advice and consent. Johnson believed that the Act was unconstitutional (and the Supreme Court, years after his Presidency, agreed in 1926), and intentionally violated it, to "test the waters." Radical Republicans used this violation as an excuse to impeach Johnson, who was acquitted by one vote in the Senate. In the election of 1868, Ulysses Grant was nominated for the Republican ticket and won on an incredibly small margin. Republicans noticed that if they did not act swiftly to protect the voting rights of blacks, they might soon lose a majority. Thus, Congress passed the Fifteenth Amendment in 1869, which enforced that the suffrage of male citizens shall not be denied on account of race. This was a major blow to the women's movement, as it was the first time gender was deliberately placed into the Constitution. Republicans claimed that if the amendment had included both race and gender discrimination clauses, it would have never had a chance to pass in Congress. African-Americans in Congress. A number of African-Americans were elected for the first time in American history during this period. With the Reconstruction Acts sending federal troops in the southern states where African-Americans held majorities in South Carolina and Mississippi, and nearly equal numbers with whites in Louisiana, Florida, Georgia, and Alabama, Blacks were elected to Congress from these states. John Willis Menard was elected in the 2nd District of Louisiana in 1868. His challenger, Caleb Hunt, filed an objection with the election result and the House of Representatives, upon hearing arguments from both candidates, decided to seat neither of them. Hiram Revels was elected by the Mississippi Senate by an 81-15 margin to finish the term of Mississippi Senator Albert G. Brown, who vacated the seat during the Civil War. Revels served from February 23, 1870 to March 3, 1871. Joseph Rainey was elected to the US House of Representatives from South Carolina's 1st District in the elections of 1870. He was the longest serving African-American member of congress prior to William L. Dawson in the 1950's. Blanche Bruce was elected to serve a full term in the US Senate by the Mississippi state senate in 1871. Bruce was the only former slave to ever serve in the US Senate. Alaskan Purchase. Beginning in the 1770's the Russian Empire began colonizing Alaska. On March 30, 1867 the American Government purchased Alaska from the Russian empire for 7.2 million dollars. The decision was widely ridiculed at the time, but the purchase was later proven to be a bargain when gold and oil were discovered there. The few Russian settlers in the territory were given three years to return to Russia, with the option of staying and becoming American citizens. The Panic of 1873. The Panic of 1873 was the first depression experienced by America and Europe following the Civil War. The depression was a result of the fall for an international demand for silver. Germany stopped using the silver standard after the Franco-Prussia war. The United States enacted the Coinage Act of 1873 which shifted the backing of our monetary system with gold and silver to just gold. The act immediately depreciated the value of silver and hurt western mining operations. Another factor that influenced the Panic of 1873 was the risky over investment into railroad companies that would not bring quick returns. The Jay Cook and Company was a United States bank that declared bankruptcy on September 18, 1873. The bank went under as a result of over investment in the railroad business. As a result, the New York Stock Exchange closed for ten days starting September 20, 1873. 89 of 364 railroad companies failed during the depression. Real estate values, wages, and profits by corporations decreased over the course of the panic. Thousands of businesses fell during the depression as well. The depression was a major highlight in President Grant's second term. The Great Railroad Strike of 1877. The strike started on July 14, 1877 in Martinsburg, West Virginia. The strike was caused by wage cuts from the Baltimore and Ohio Railroad Company. The workers refused to let the railroad operate. State militia was sent in to quell the strike but would not fire upon the strikers. Governor Henry Mathews called upon federal troops to put down the strike and resume operations of the railroad. The strike spread to Cumberland, Maryland. Troops in Maryland fired upon the mob of strikers and killed ten rioters. The strike occurred in Philadelphia, Baltimore, Pittsburgh, and even spread to St. Louis. The strikes resulted in millions of dollars of property damage the casualties of many. The great strike lasted 45 days, after finally being put down by federal troops from city to city. Republicans fall from power. Grant's presidency would bring about the decline of the Republican Party. He appointed a great number of corrupt officials to federal positions and to his cabinet. Many split with the party over that issue. Others grew tired of Reconstruction and proposed reconciliation with the South in a peaceful manner. These people called themselves Liberal Republicans, and nominated Horace Greeley to run against Grant in 1872. The Democrats also endorsed Greeley. Despite wide support, Grant won the election of 1872 decisively. During the election season, Liberal Republicans were busy pushing the Amnesty Act through Congress, and in May 1872, it passed. The Amnesty Act pardoned most former Confederate citizens, and allowed them to run for office. The act restored the rights to the Democratic majorities in the South. Soon, Democrats had control of the Virginia and North Carolina governments. In states with black Republican majorities, the Ku Klux Klan (formed after the civil war as a white supremacist group) terrorized Republicans and forced them to vote Democratic or not at all. By 1876, Republicans controlled only three states in the South: Florida, Louisiana, and South Carolina-- all of which were still occupied by Union troops. Republicans continued to decline during Grant's second term, after many high level political scandals came to light. Most shocking to the public was that a scandal involved the Vice President, and another involved the Secretary of War. The Northern population's confidence in the party was shaken even more when the nation slipped into a Depression that same year. In the congressional elections of 1874, Republicans would suffer huge losses in both houses, and for the first time since before the start of the Civil War, Democrats were able to gain control of a part of Congress (the House). Congress no longer was able to be committed strongly to Reconstruction. In the election of 1876, Democrats nominated New York governor S.J. Tilden to run, and the Republicans nominated Ohio governor Rutherford B. Hayes. On election day, it seemed that Tilden would win by more than 250,000 votes. But the seven, four, and eight electoral votes from South Carolina, Florida, and Louisiana, respectively, were disputed (Northern troops still occupied these states). Also, one of Oregon's three electoral votes was disputed. If Hayes won all 20 votes, he would win the election. Congress created a special commission of seven Democrats, seven Republicans, and one independent to review the election and decide a winner. But the independent resigned, and a Republican was appointed to take his place. The commission voted along party lines to award Hayes the election, but Democrats warned that they would fight the decision. Republican and Democratic leaders secretly met up to draw up a compromise, and the result of the meeting was the Compromise of 1877. Proclaiming that Hayes would win the election, troops left the South and more aid was given to the South; it marked the end of Reconstruction. Ultimately, Reconstruction and the Compromise itself would be failures, as Democrats refused to hold up their end of the compromise, which was to protect the rights of African Americans in the South. The period after Reconstruction saw the rise of the Democratic "Redeemers" in the South. The Redeemers vowed to take back the South from Republican rule, which had been ousted after the 1876 election. They passed Jim Crow laws, which segregated blacks and whites, and put voting restrictions on blacks that wouldn't be outlawed until the next century. Jim Crow laws were challenged in "Plessy v. Ferguson", when the Supreme Court voted to uphold the laws if and only if segregated facilities remained "separate but equal." Sinmiyangyo. The United States expedition to Korea in 1871, also known as Sinmiyangyo (Western Disturbance of the Year Sinmi year) was the first American military action in Korea. It took place predominantly on and around the Korean island of Ganghwa. America sent a military expeditionary force to Korea to support an American diplomatic delegation sent to establish trade and diplomatic relations with Korea, to ascertain the fate of the General Sherman merchant ship, and to establish a treaty assuring aid for shipwrecked sailors. The isolationist Joseon Dynasty government and the assertive Americans led to a misunderstanding between the two parties that changed a diplomatic expedition into an armed conflict. The United States won a minor military victory, but the Koreans refused to open up the country to them. As the U.S. forces in Korea did not have the authority or strength to press the issue, the United States failed to secure their diplomatic objectives. Religion During the 19th Century. The New Wave of Jewish Immigration. Between the years of 1820 and 1880, about 250,000 Jews came to the U.S. The immigration was not only based in a troubled European economy but also in the 1848 failure of liberal revolutions in the German states. Railways and the great steamer ships opened up immigration to America in the later 1800s and early 1900s. Yet the trip was dirty, uncomfortable, and dangerous. While prosperous families in upper-class compartments had private cabins, most of the immigrants went across in steerage class: three hundred tightly-packed men, women and children, sleeping on double- or even triple-bunk beds. The beds were about six feet long and two feet wide, with only two-and-a-half feet separating each bunk. Their goods, such as they were, rested on the bunk. The people had little or no running water, the stench was incredible, and there was little recourse against vermin. One witness, a Sophia Kreitzberg, said, "When you scratched your head [. . .] you got lice on your hands." This close environment joined with the background of privation and the strain of long travel to cause break-outs of diseases. The Jews were served unkosher meat and soup, which many refused to eat. Instead they nibbled on what they had brought with them, mostly dried fruit, hard bread, or stale cheese. At the other end of the passage was often a few ports of entry to the United States. From 1855 to 1890 this was predominantly Castle Garden (also called Castle Clinton), a point at the tip of Manhattan Island in what would later be called New York City. From 1890 onward a vast number of immigrants came through the new reception area at Ellis Island. In these places human beings were duly categorized by ship, country of origin, past employment, and other relevant information. (This next footnote shows one such document from Castle Garden.) Catholic America. Catholic Immigration. A massive influx of immigrants from Ireland, Italy, Germany, Austria-Hungary, and the Russian Empire (most notably Poland) caused a dramatic increase in the U.S. Catholic population during the latter half of the 19th century. From 1840 to 1851 Ireland suffered famine and oppression. The other immigration was caused by nationalization and national upheavals. By 1850 Catholicism had become the United States' largest religious denomination: between 1860 and 1890 the Catholic population tripled, mostly because of immigration. This massive influx of Catholics to the United States eventually led to a significant increase of power for the Catholic church. Persecution. American Protestants often feared the increasing power and prominence of American Catholics. In pre-Civil War America anti-Catholic prejudice was shown through the "Know Nothings" and the American branch of the Orange Institution. After the war the American Protective Association, and the Ku Klux Klan regularly persecuted and discriminated against Catholics with such acts as The Philadelphia Nativist Riot, "Bloody Monday", and the Orange Riots of New York City in 1871 and 1872. Nativism. The severe anti-Catholic activities revealed the sentiments of Nativism, which encouraged all "native-born American men" to rise up against foreigners. (This appeal went out to European-descended Protestants, of course, rather than the actual Native Americans who lived there before the had settled here.) The first Nativist publication was called "The Protestant"; its first edition sold on January 2nd, 1830. Its editor was George Bourne, and as he wrote, "the goal of the paper is centered around the denunciation of the Catholic faith" Anti-Catholic rhetoric was occasionally met with violence; yet the Nativists produced one of the greatest violent acts of the 1830's. On August 10, 1834, forty to fifty people gathered outside the Ursuline Convent school and burned it to the ground. In 1834 F. B. Morse, a nativist leader who was a professor of sculpture and painting at New York University, wrote "The Foreign Conspiracies Against the Liberties of the United States", in which his basic message is centered around protecting the American birth right of liberty. The concern, and fear of the foreign and Catholic communities grew out of the Protestant fear of the monarchial tendencies of Catholicism, during this time urban areas were also starting to grow rapidly with the massive influx of immigrants who all congregated and lived in the same areas. Nativists saw this as an act of "clannishness", and an attempt to avoid or resist "Americanization." With the success of Morse, and his contemporary Lyman Beecher, the nativist movement reached a point where the public did not care whether the stories they heard were true or false, but began to accept works of fiction as truth as well. In 1836 Maria Monk authored a worked called "Awful Disclosures of the Hotel Dieu Nunnery of Montreal." In her book she tells of her experiences with Catholicism, which involved forced sexual intercourse with priests and the murdering of nuns and children, the book concludes with her [Maria] escaping to save her unborn child. Monk's mother denies her work, and said that Maria was never in a nunnery, and that a brain injury Maria received as a child may have been the cause of her stories. In the Midwest and northern sections of the country Catholics were seen as incapable of free thought and were said to be "anti-American Papists" because it was thought that they took every direction from the Pope in Rome. During the Mexican-American war Mexican Catholics were displayed in the media as silly or stupid due to their "Papist superstition". It was because of the general attitude in America about Catholics that about 100 American Catholics, mostly recently immigrated Irish, fought against the United States in the Mexican-American war. These men fought for the Mexicans and were known as "Saint Patrick's Battalion (). In 1850, Franklin Pierce presented several resolutions that would remove the restrictions on Catholics from holding public office in New Hampshire, these resolutions that were, at the time, considered "pro-Catholic' were defeated (Battle of Religious Tolerance," The World Almanac, 1950, 53). However as the 19th century passed, hostilities between Catholics and Protestants eased due to the fact that many Irish Catholics fought alongside Protestants during the Civil War, for both the North and the South. Education. Ex-slaves everywhere across the nation reached out for education. Blacks of all ages really wanted to know what was in the books that had been only permitted to whites. With freedom they started their own schools and the classes were packed days and nights. They sat on log seats or the dirt floors. They would study their letters in old almanacs and in discarded dictionaries. Because the desire to escape slavery's ignorance was so great, ignoring their poverty, many blacks would pay tuition, sometimes $1 or $1.50 a month. Blacks and their white allies also saw a need for colleges and universities, in this case to train teachers, ministers, and professionals for leadership. There were seven colleges founded by the American Missionary Association, Fisk and Atlanta Universities, between 1866 and 1869. The Freedoms Bureau helped to establish Howard University in Washington D.C. As well as Northern religious groups, such as the Methodists, Baptists, and Congregationalists, supported dozens of seminaries and teachers’ colleges. The earliest forms of education that blacks received was from the missionaries to convert them to Christianity. The education of blacks was very low during the civil war, until Lincoln issued the Emancipation Proclamation in 1863. The Department of Education was developed in 1867 to help start more effective schools systems. Howard University was developed in Washington D.C. for black youth “in the liberal arts and sciences.” The first public school day was in Boston in 1869. Technology. Just before the Civil War a vast source of petroleum was discovered in and around Titusville, Pennsylvania, and after the War this began to be exploited. At first oil was used for medicinal purposes alone. However, as the supply increased, it also began to be used for industrial purposes, and instead of whale oil. The dangerous and expensive whaling industry collapsed. While some cities used coal gas for night illumination, others began to use oil lamps, and major cities were lit at night. Petroleum, lamp oil (for the great engine lamp), and machine oil increased the usefulness of railways. Information could be transmitted across great distances via the telegraph. In the 1870s and 1880s inventors vied to transmit a human voice. The two major competitors were the Scots-born Alexander Graham Bell and Elisha Gray. In the year 1875, Alexander Graham Bell used an electromagnetic machine to transmit the sound of a steel reed. On February 14, 1876, a partner of Bell presented his patent to the patent office in Washington, D.C., on the same day as his rival Elisha Gray. Three weeks later, on March Seventh, Bell’s patent won out and was granted. Native Americans After The War. The injustices that Native Americans dealt with during the Civil War did not go away at War's end. The U. S. National government made it clear that though these people were indigenous to the continent, they were not going to be citizens of the country. The Native Americans were forced to live out in the west on reservations. Their travel was restricted and scrutinized by government agents who monitored them. Traveling off the reservations to hunt, fish or even visit the neighboring reservations was frowned upon by the Bureau of Indian Affairs. Subsequently, the Bureau instituted a pass system in order keep them under control. This system required the Natives to get a pass from the agents before they were allowed off the reservation. White settlers also took issue with Indians traveling on trains. However, the Central Pacific Railroad in Nevada granted the Native Americans permission to ride on top of the trains in exchange for their railroads being allowed to cross through the reservations. Many Indian agents were unhappy with this free travel arrangement. They began writing letters to the Bureau to stop it. One of the agents commented that "The injurious effects of this freedom from restraint, and continual change of place, on the Indian, can not be overestimated." With the 14th amendment the civil rights acts were contrived. For the Indians however, their positioning was made clear. The Civil Rights Act of 1866 states “That all persons born in the United States, and not subject to any foreign power, excluding Indians not taxed, are hereby declared to be citizens of the United States." Battle at Little Bighorn. In 1876, after a few uneventful confrontations, Col. George A. Custer and his small cavalry came across the Sioux and some of their allies at the Little Bighorn River. To force the large Indian army back to the reservations, the Army dispatched three columns to attack. One of these groups contained Lt. Custer and the Seventh Cavalry. They spotted the Sioux village about fifteen miles away just along the Rosebud River, Custer also found a nearby group of about forty men. He ignored orders to wait, and decided to attack before they could alert the main party. He was unaware of how much he was outnumbered. The Sioux and their allies had three times as much force. Custer divided his forces in three, He sent troops under control of Captain Frederick Benteen to try to stop them from escaping through the upper valley of the Little Bighorn River. Major Marcus Reno job was to pursue the group, then cross the river, and attack the Indian village in a conjunction with the remaining troops under his command. He Intended to strike the Indian camp from the north and south, but he had no idea that he would have to cross a rough terrain in order to achieve this. As the Indians descended Custer ordered his men to shoot their horses and stack the carcasses in front of them in order to form a wall. But this did not protect them against bullets. In less than an hour, Custer and all his men were killed in one of the worst American military disasters of all time. After one more day of fighting, Reno and Benteen's now unified forces fled when the Indians stopped fighting. They [who?] knew two more columns of soldiers were coming towards them, so they escaped toward them. The massacre of Custer's final battle eclipsed any success he had had in the Civil War. Custer was defeated and killed at the Battle of the Little Bighorn on June 25, 1876, while fighting Native American tribes in a battle that has come to be known as "Custer's Last Stand". Women's History of the Period. Victoria Woodhull. In 1872 Victoria Woodhull became the first woman to run for President of the United States. She was nominated by the Equal Rights Party on May 10. Though it is undisputed that she was the first female to run for president, the legality of her petition is questioned; her name didn't actually appear on the ballot and she was under the age of 35 which is the required age for a presidential candidate according to the constitution. Woodhull did not receive any electoral votes, but evidence supports that she received popular votes that were never counted. Woman's Christian Temperance Union. The Woman's Christian Temperance Union was formed on December 22, 1873. Fredonia, New York is credited as being the birthplace of the group. The temperance movement was a social movement that pushed for the reduction of alcohol consumption. The movement spread all over the country, and women would go to bars and drug stores to sing and pray. The National Woman's Christian Temperance Union was established in 1874 in Cleveland, Ohio. The women demonstrated use of non violent protestation of the consumption of alcohol by praying in saloons. Often, they were denied entrance and yelled at by patrons. The movement ultimately contributed to prohibition in America's future. 

A student decided to silver plate his/her mother's forks. S/he set up the apparatus above. S/he set the power pack and variable resistor so that the ammeter read 0.1 A. S/he allowed the fork to remain in the electrolysis apparatus for 10 min. Once s/he removed the fork s/he tested it by trying to scratch the silver off. S/he found that the layer of silver was too thin and decided that s/he would have to take steps to make the layer a lot thicker. Which one of the following steps would not produce a thicker layer of silver ? 

This is the correct answer. Well done! Increasing the resistance would "decrease" the current and so decrease the thickness of the silver. On to the next question» 

I'm sorry but this is the wrong answer. Increasing the voltage would increase the current making the layer of silver thicker. But the question asked which one of the four possible measures would "not" increase the thickness of the layer. «Back 

I'm sorry but this is the wrong answer. Increasing the time gives the silver more time to build up. This means the layer "will" be thicker. «Back 

I'm sorry but this is the wrong answer. Increasing the concentration of silver nitrate will mean there are more ions to carry the current. The current will therefore be larger and so the silver layer will be "thicker". The question asked which method "'would not" produce a thicker layer. «back 





The is a worldwide collection of computer networks that began as a single network that was originally created in 1969 by ARPA (Advanced Research Projects Agency), a U.S. government agency that was far more interested in creating projects that would survive a nuclear war than in creating anything useful for the civilian population. In its original form, ARPANET, the U.S. government hoped to create a network of computers that would allow communication between government agencies and certain educational centers that would be able to survive a nuclear explosion. It is doubtful that the original founders of ARPANET foresaw what we now know as "the Internet." From its humble beginnings as a military project, the ARPANET grew slowly throughout the 70's and 80's as a community of academics accomplished the truly monumental task of hammering out the building blocks of this new, open, modular conglomeration of networks. In addition to the U.S. ARPANET, other countries also developed their own computer networks which quickly linked up to ARPANET, such as the UK's (1983 onwards), and Australia's (mid-1970s until replaced). Connecting these together would help develop a global internetwork. The various protocols, including IP, TCP, DNS, POP, and SMTP, took shape over the years, and by the time the World Wide Web ( and HTTP) was created in the early 90's, this "Internet" had become a fully functional, fairly robust system of network communication, able to support this new pair of protocols which eventually turned the Internet into a household word. While a large portion of users today confuse the Web with the Internet itself, it must be emphasized that the Web is only one type of Internet application, and one set of protocols among a great many which were in use for over a decade before the Web entered into the public awareness. The Web is a subset of the Net. Email is not a part of the Web, and neither are newsgroups, although Web designers have developed web sites through which users, the world over, commonly access both of these much older forms of Internet media. While the Net is a largely abstract phenomenon, it cannot (at least, not yet) be accurately equated with the concept of "cyberspace" as depicted in science fiction. If "judgement day" were to occur as depicted in the latest "Terminator" film, much of the Internet would survive it, but most of the electrical and data infrastructure by which we access the net would not. The line which currently demarcates the "digital divide" would shift dramatically to a point where it would leave only a small segment of humanity in virtual touch. This limitation, however, will slowly be overcome as wireless technologies continue to proliferate and wired technologies become increasingly cheaper. In March 1972 ARPA became known as DARPA, the Defense Advanced Research Project Agency, and then went back to ARPA in February 1993 and back to DARPA in March 1996 and has been ever since. It was originally created as ARPA in 1958 in response to the launching of Sputnik. The launch of Sputnik made America realize that the Soviet Union could exploit military technology. DARPA has contributed to the creation of ARPANET as well as the Packet Radio Network, the Packet Satellite Network and the Internet. As well as research into the Artificial Intelligence field commonly referred to as AI. By the late 1970's the Department of Defense had adopted BSD UNIX as the primary operating system for DARPA. 

Foreword. There are many books in the market nowadays that have written about the puzzles BUT most of the books only dealt with narrow scope of that particular puzzles only. For this Wikibook, it is the intention of the contributors to give the readers some tidbits/history and examples related to that puzzles mentioned and of course there are some puzzles which provided entertainment / interactive question for readers themselves at the end of the page. If you got stuck after trying, it wouldn't hurt to take a peek on how to solve the puzzles itself (after all puzzles is supposed to be fun). Hence all said, please enjoy these puzzle collection, avoid any frustration, and fight boredom. Puzzle resources.  __NOEDITSECTION__ 

The peoples of Europe have had a tremendous effect on the development of the United States throughout the course of U.S. history. Europeans "discovered" and colonized the North American continent and, even after they lost political control over its territory, their influence has predominated due to a common language, social ideals, and culture. Therefore, when endeavoring to understand the history of the United States, it is helpful to briefly describe their European origin. Greece and Rome. Ancient Greece. The first significant civilizations of Europe formed in the second millennium BCE. By 800 BCE, various Greek city-states, sharing a language and a culture based on slavery, pioneered novel political cultures. In the Greek city of Athens, by about 500 BCE, the male citizens who owned land began to elect their leaders. These elections by the minority of a minority represent the first democracy in the world. Other states in Greece experimented with other forms of rule, as in the totalitarian state of Sparta. These polities existed side by side, sometimes warring with each other, at one time combining against an invading army from Persia. Ancient Athens is known for its literary achievements in drama, history, and personal narrative. The individual city-states did not usually see themselves as a single entity. (The conqueror "Alexander the Great", who called himself a Greek, actually was a native of the non-Greek state of Macedon.) The city-states of Greece became provinces of the Roman Empire in 27 BCE. Rome. The city of Rome was founded (traditionally in the year 753 BCE). Slowly, Rome grew from a kingdom to a republic to a vast empire, which, at various points, included most of present-day Britain (a large part of Scotland never belonged to the empire), France (then known as Gaul), Spain, Portugal, Italy, Greece, Turkey, Iraq, Palestine (including the territory claimed today by the modern state of Israel), Northern Arabia, Egypt, the Balkans, and the entire north coast of Africa. This empire was maintained through free-born or adoptive citizenship, citizen education and indoctrination, a large and well-drilled army, and taxes directed by a large bureaucracy directed by the emperor. As each province produced more Roman citizens, the state became hard to maintain. Whole kingdoms in the north and east, and the invading peoples we know as the Germanic tribes (the Ostrogoths and Visigoths and the Franks) sat apart from the system. After the death of one emperor in 180 CE, power struggles between the army and a succession of rulers of contested origins produced anarchy. Diocletian (243 - 316) reinstated the Empire by 284. Rome regained territory until 395, when the Empire was so large that it had to be divided into two parts, each with a separate ruler. The two halves sat uneasily together. The East, which considered itself the heir of Alexander the Great, spoke Greek or a dialect, while the West spoke Latin. The Eastern Empire survived until 1453, but the system to maintain the Western Empire broke apart. Plagues and crop failures troubled the world. In 476, Germanic tribes deposed the boy who was then the Emperor. Roman roads fell into disrepair, and travel became difficult. Some memories remained in the lands which had once known Roman rule. The supreme rulers of various tribes called themselves "king," a distortion of the Roman word Caesar. The Roman Empire to the Holy Roman Empire. After Rome's fall, monks from Ireland (which had never known Rome) spread Catholic Christianity and the culture and language of the Western Roman Empire across Europe. Catholicism eventually spread through England (where the Germanic tribes of the Angles and the Saxons now lived)and to the lands of the post-Roman Germanic tribes. Among those tribes, the Franks rose to prominence. "Charlemagne" (742 - 814), the King of the Franks, conquered great portions of Europe. He eventually took control of Rome. The senate and the political organs of Rome had disappeared, and Charlemagne did not pretend to become the head of the Church. Charlemagne's domain, a confederation of what had been Roman Gaul with Germanic states, was much smaller than Diocletian had known. But prestige came with identity with the past, and so this trunk of lands became The Holy Roman Empire. Charlemagne's descendants, as well as local rulers, took their sanction from the Church, while the Church's pope influenced both religious and political matters. The result of political stability was technological advance. After the year 1000, Western Europe caught some of the East's discoveries, and invented others. In addition to vellum, Europeans now started making paper of rags or wood pulp. They also adopted the wind and water mill, the horse collar (for plows and for heavy weights), the moldboard plow, and other agricultural and technological advances. Towns came into being, and then walled cities. More people survived, and the knights and kings over them grew restive. Viking Exploration of North America. In the eighth century, pushed from their homes in Scandinavia by war and population expansion, Norsemen, or Vikings, began settling parts of the Faeroe, Shetland, and Orkney Islands in the North Atlantic. They went where ever treasure was, trading as far as Byzantium and Kiev in the East. In the West they raided from Ireland and England down to the Italian peninsula, sailing into a port, seizing its gold, and murdering or enslaving its people before fleeing. They began settling Iceland in approximately 874 CE. A Viking called "Erik the Red" was accused of murder and banished from his native Iceland in about 982. Eric explored and later founded a settlement in a snowy western island. Knowing that this bleak land would need many people to prosper, Eric returned to Iceland after his exile had passed and coined the word "Greenland" to appeal to the overpopulated and treeless settlement of Iceland. Eric returned to Greenland in 985 and established two colonies with a population of nearly 5000. "Leif Erikson", son of Erik the Red, and other members of his family began exploration of the North American coast in 986. He landed in three places, in the third establishing a small settlement called Vinland. The location of Vinland is uncertain, but an archeological site on the northern tip of Newfoundland, Canada (L'anse aux Meadows) has been identified as the site of a modest Viking settlement and is the oldest confirmed presence of Europeans in North America. The site contains the remains of eight Norse buildings, as well as a modern reproduction of a Norse longhouse. But the settlement in Vinland was abandoned in struggles between the Vikings and the native inhabitants, whom the settlers called "Skraelingar". Bickering also broke out among the Norseman themselves. The settlement lasted less than two years. The Vikings would make brief excursions to North America for the next 200 years, though another attempt at colonization was soon thwarted. By the thirteenth century, Iceland and Greenland had also entered a period of decline during the "Little Ice Age." Knowledge of their exploration, in the days before the printing press, was ignored in most of Europe. Yet the Vikings are now considered the true European discoverers of North America. The influence of their people outlasted even the terrible raids, and their grandchildren became kings and queens. For example, a branch of Viking descendants living in France, the Normans, conquered England from the Anglo-Saxons in 1066. The Great Famine and Black Death. The Little Ice Age led to European famines in the years 1315-1317 and in 1321. In the year 1318 sheep and cattle began to die of a contagious disease. Farmers could not support the growing population. And then, in 1347, some Genoese trading ships inadvertently brought a new, invasive species of rat to Europe. These rats carried bubonic plague. Plague was also called the Black Death, from the darkened skin left after death and from its deadly reach. It had three strands: "bubonic", "pneumonic", and "septicemic". In bubonic plague fleas carried by the rats would leave their hosts and bite people. The masses of bacteria would flow through the human system, killing cells and leaving their refuse in lymph nodes in the armpit, groin, and neck. These nodes would swell and turn black, creating bubos. Infection could also spread into the lungs, so that a person might cough or sneeze the germ into the air. This created pneumonic plague, spreading disease into spaces where people gathered and where rats dared not go. It also spread through contamination of food. The last form of disease, and its most deadly, was septicemic. This attacked the blood, leaving stretches of pale skin looking black, and killing the person within hours. Surviving laws of cities and guilds regulate public cleanliness and penalize adulteration of food. They cannot show how strictly these laws were applied. And they show no knowledge, of course, of germ theory and the need for sterilization. Older systems such as the few public baths which remained from the days of colonial Rome were seen as sinful and dangerous, invitations to the plague. The dwelling places of survivors of pre-Christian Rome, the Jews who were forced to live apart, were attacked by mobs who attributed the Black Death to their poisoning Christians' wells. The responses to plague can be seen in the records left by survivors -- one third of the population of Europe died in repeated waves of disease -- and in the subsequent changes in society. Airplanes and satellites show the foundations of plague-era towns which were emptied by the disease. In just one square mile of pre-plague Europe there are reports of there being 50,000 people. In large cities, families would flee or lock themselves away, trying to keep themselves from death. Other families were locked in by city authorities. This is the beginning of the modern system of quarantine. Some branches of the family would not be among those so helped. The Black Death seemed erratic, sometimes taking people deemed good and pious, sometimes not. One priest or church prelate might die, and another survive. And a living priest might give no aid to other survivors. Some critiques of the Church which had become spread through most of Europe date from this era. Although some lands became waste through lack of tilling, those people who survived grasped the property of those who had not. Europe then had a land-based economic system. Rich people became richer. There began a labor shortage: the farmers who survived needed hands to take in their crops. The wages of farm hands began to rise. In the surviving towns they needed people to guard the gates: in the courts they began to look for rising young men. Cities became more powerful in the depleted lands, and authority grew more centralized. Education. During the Middle Ages Western society and education were heavily shaped by Christianity as expressed through the Roman Catholic Church. Towns, courts, and feudal manors had their priests, monasteries and nunneries had their "scriptora" or libraries, and after the 11th century CE, a few cities had Universities, schools to educate men to be high-ranking clerics, lawyers, or doctors. Where children had schools, their parents paid a fee so that they might learn Latin, the language St. Jerome had used for his translation of the Bible. Latin was the language of the Church. It was also learned, along with military tactics and the rules of chivalry, by men who trained to be knights. A smattering of Latin was necessary, along with Math, even for the elementary schools which sprang up in some cities. There both boys and girls were taught literacy and math, prerequisite for acceptance as an apprentice in many Guilds. Latin across Europe created a European-wide culture: a doctor from Padua could talk to his fellow from Oxford in Latin. As in the Greek and Roman eras, only a minority of people went to school. There were not enough books, little travel, and no means of spreading standardized education. Schools were attended first by persons planning to enter religious life. Occasionally a cleric would reach out to educate a very bright peasant boy. This was one of the few ways peasants could rise in the world. But the vast majority of people were serfs who served as agricultural workers on the estates of feudal lords. They were, in effect, tied to the fate of the land. From the serf up to some high princes, the vast majority of people did not attend school, and were generally illiterate. In the rise of the Universities in the 11th century, the Church translated several manuscripts of the ancient Greek writer Aristotle into Latin. From Aristotle's emphasis upon human reason, philosophy and science, and the Church's emphasis upon revelation and the teachings of Christ, medieval scholars developed "Scholasticism". This was an philosophical and educational movement which attempted to integrate into an ordered system both the natural wisdom of Greece and Rome and the religious wisdom of Christianity. It was dominant in the medieval Christian schools and universities of Europe from about the middle of the 11th century to about the middle of the 15th century, though its influence continued in successive centuries. The idea of Man in the middle of ordered nature, and yet dominant over nature, bore fruit in the observation of natural phenomena, the beginnings of what the Western World knows as science. It also led to the exultation of system over observation, and the persistence of the Ptolemaic theories such as geocentrism among formally educated people. Noble girls were sometimes sent to live in nunneries in order to receive basic schooling. Nuns would teach them to read and write and the chores necessary to run their establishments, including spinning and weaving. (Cloth-making was a major national industry in the Middle Ages.) They taught them their manners and their prayers. Some of these girls later became nuns themselves. Christianity, Islam, Judaism. During the centuries after the fall of Rome, various flavors of Christian churches spread from Northern Africa and Armenia westward. This changed after "Mohammed" established Islam in 610 CE. Like Christianity, it spread through conversion and conflict. At its height it was also a faith of Europe, from Spain to Albania and Bosnia and their sister states. Both Prince and Caliph held that their state must have one faith, and no other belief was encouraged. When Jerusalem was reconquered by the Seljuq Turks, Christians were no longer able to go on religious pilgrimages to the Holy City. At the end of the eleventh century, "Pope Urban II" inaugurated the Crusades, urging Western European kings and great nobles to begin what would be a century and a half of warfare. Christian armies fought first to reconquer and then to hold part of the Kingdom of Jerusalem. The Crusaders ultimately failed in the face of resurgent Muslim forces. Western Europeans within Church and State argued for and against the Crusades. Despite the failure of the Crusades, militant Western Christianity persisted in Spain in an effort known as "the Reconquista" (the "reconquest"), which purged the land of the Muslims who had arrived there in 711. By the fifteenth century, the Muslims were confined to the kingdom of Granada, which bordered the Mediterranean Sea in the southern side of the Iberian Peninsula. Granada finally fell in 1492 to the Spanish Christians, ending the reconquista. Rome had destroyed Israel in 70 CE, but allowed a remnant of her people to survive. Fortified by rabbinic culture and centered on the Torah, they became a resilient group. They survived as the known world became Christian. They spread, as traders, through East and West. However, Christian relationships with Jews were punctuated by hatred. They were held to be guilty of Jesus' death, and they were supposed to be evil because they had not converted to Christianity. They were forced into the notorious ghettos, usually built on waste or undesirable land. They were forced to wear strange clothing which marked them off, and to pay heavy taxes. Christians spread rumors that Jewish officials sometimes kidnapped and killed young boys for their sacrifices. Sometimes a mob might break into the ghetto, killing some people. Individuals were sometimes forced to convert. One of the effects of European nationalism was the expulsion of a country's Jewish population. England was the first, in the 13th century; later, a reviving France; later still, Spain and Portugal. Some men crewing the ships in the Age of Exploration were Jews, often practicing their faith in hiding. For a period in the late Twelfth century there two sets of Popes, a line in Rome and another in the French city of Avignon under the sway of that Court. In reaction against this, the Church centralized its powers in parallel to what nations were doing. There had been dissension before, in medieval England's Lollards and later with the Czech priest John Hus. However, it was only in an age after the printing press, when people began printing the Bible in their own languages, when "Martin Luther" founded the Protestant church. England's "King Henry VIII", who had won the title "Protector of the Faith" for a work defending the Pope, later left it to become the head of the Church of England. This division of Western Christianity created religious minorities, who were persecuted throughout Europe. Among these were the Pilgrims, who helped settle America. The Renaissance. Another, more humble result of the plague was the accumulation of rags left over from clothing. These were quickly used to make paper. Books were very rare during the Middle Ages, and the monks who made them chained them to their shelves. It took one year for a man to make one book. In that climate Bibles took priority: we have only one copy of Democritus's most famous work surviving from this period. As rag paper replaced velum, books began to become more plentiful. The supply was augmented by Europe's adoption of "Johannes Gutenberg" 's fixed-type printing press in the 15th century. There had been earlier leaps in culture, including the wave of population and technological adaptation in the 12th century. This left its mark in increased population and the Roman Catholic Church's adoption of Aristotle. Yet the press made it possible for knowledge to have a foothold in society. Inventions in one place could be explained and adopted across a continent. The Greeks fleeing the fall of Byzantium brought their knowledge of ancient Greek culture to the West. The Bible, the basic book of Christendom, could be pored over by laypeople, and reading it could be learned by more people than ever before. Learning was no longer solely the province of the Church. If the Bible was first off of the presses, pseudo-science and science followed soon after. The European witch trials were one result of the new medium. The questioning of scientific consensus was another. "Andreas Vesalius" published his observations about the circulatory system. Books discussing the theories of "Nicolaus Copernicus" and "Galileo Galilei" demolished the old geocentric theory of the cosmos. The arts were not neglected. Giorgio Vasari's biographical Lives introduced such new artists as "Leonardo Da Vinci" and "Michelangelo Buonarroti" to the larger world. A later time called this growth of knowledge the Renaissance. They said it began in the Italian city-states, spreading throughout most of Europe. The Italian city of Florence was called the birthplace of this intellectual movement. Books spread the Crusader's newly found experience and knowledge of the Mediterranean, a region whose technology was at that time superior to that of western Europe. Books written about traders, adventurers, and scholars spread knowledge of Chinese technology such as gunpowder and silk. They spread writings of the ancient world which had been lost to Europe, and nurtured a taste for new foods and flavors. They spread pictures of ancient Greek statues, Moorish carpets, and strange practices. In the fifteenth century, the Mediterranean was a vigorous trading area. European ships brought in grains and salts for preserving fish, Chinese silks, Indian cotton, precious stones, and above all, spices. White cane sugar could be used to preserve fruit and to flavor medicine. Cinnamon was medicine against bad humors as well as preservative and flavoring, part of the mysterious "poudre douce", and now available even to some European common people. Toward Nation-States. The great age of exploration was undertaken by nation-states, cohesive entities with big treasuries who tended to use colonization as a national necessity. When such nations as Portugal, England, Spain, and France became stronger, they began building ships. The Rise of Portugal. The Italian peninsula dominated the world because of its position in the Mediterranean Sea. Universities in Padua, Rome, and elsewhere taught men from East and West. Above all, principalities such as the Republic of Venice and Florence controlled trade. Genoa and Venice in particular ballooned into massive trading cities. Yet there was at yet no nation of Italy, so each city's riches belonged only to that city. Individual cities used their monopoly to raise the price of goods, which would have been expensive in any case, because they were often brought overland from Asia to ports on the eastern Mediterranean. The mad prices, in turn, increased the desire of purchasers to find other suppliers, and of potential suppliers to find a better and cheaper route to Asia. Portugal was just one of many potential suppliers, with a location which extended its influence into the Atlantic and down, south and east, to Africa. Prince Henry, son of King John I, promoted the exploration of new routes to the East. He planned Portugal's 1415 capture of Ceuta in Muslim North Africa. He also sponsored voyages that pushed even farther down the West African coast. By the time of his death in 1460, these voyages had reached south to Sierra Leone. Under King John II, who ruled from 1481 to 1495, Bartolomeu Dias finally sailed around the southernmost point of Africa to the Cape of Good Hope (1487-1488). In 1497-1499 Vasco da Gama of Portugal sailed up the east coast of Africa to India. The Portuguese colonized and settled such islands in the Atlantic as Madeira, Cape Verde, the Azores, and Sao Tome. These islands supplied them with sugar and gave them territorial control of the Atlantic. West Africa was more promising, not only unearthing a valuable trade route to India, but also providing the Portuguese with ivory, fur, oils, black pepper, gold dust, and a supply of dark skinned slaves who were used as domestic servants, artisans, and market or transportation workers in Lisbon. They were later used as laborers on sugar plantations on the Atlantic. France, England, And The Hundred Years' War. King Harold of England faced William the Norman usurper after defeating the last Viking forces holding the North of England. And when William the Norman became William the Conqueror, he held England by consolidating the nation. His army was ruled from newly built castles, and had the best technology of the time. His sons and their sons fought the original inhabitants of their country in Scotland and Wales, pushing their boarders. They also claimed the right to their ancestral Normandy, in what is now France. (There had long been rivalry between England and France over the wool trade.) By a English king's marriage to an former queen of France, Eleanor of Aquitaine, the king subsequently claimed Aquitaine. During the years of 1337 to 1453, the kings of England claimed the whole of France and beyond, fighting The Hundred Years' War there and in the Low Countries. For some years the English threatened Paris, and there was a question whether the small area of France proper would be entirely conquered. The early stages of the war marked by English victory against a demoralized French people and their Prince. But around 1428, a young peasant girl from Lorriene, France named Joan of Arc approached a garrison of the French army. She told them that Saint Michael, Saint Catherine, and Saint Margaret had told her to lead the army to victory. She said that God had come to her in a dream, and told her how to defeat the English. She claimed that God said Prince Charles of France needed to be crowned in order for France to claim victory over the English. After gaining the new French king, support of the populace, and several key victories, Joan was sold to the English for treason and as an appeasement. Yet France had been renewed. It pushed back against the invaders, and took back most of the land. In the next few centuries it was to remove England from Continental Europe. The Hundred Years' War devastated both countries. But it ultimately turned both of them into stronger, colonial powers. The Hundred Years' War created opportunities for wealth and advancement for the knights of both countries. The Chivalric code showed great influence during this period. France absorbed Aquitaine, Castile, and Normandy itself, prosperous areas. The twin strokes of the plague and the Inquisition weakened opposition to the French king's rule. And English centralization continued as its own royalty sought service of serf and Baron. A group of new dialects, Middle English, came out of the tug between Norman lord and Anglo-Saxon peasant. If the nation could not get new land in Europe, it could use its ships to sail elsewhere. Review Questions. 1. Explain how one of these late medieval devices affected prosperity: the wind mill; the horse collar; the printing press. 2. How did the plague infect individuals? How did mass death affect society? 3. What was the importance of these four men to the Renaissance? Galileo Galilei, Nicolaus Copernicus, Andreas Vesalius, Leonardo Da Vinci. 

Chemical equations are a convenient, standardized system for describing chemical reactions. They contain the following information. The final two points are optional and sometimes omitted. Anatomy of an Equation. formula_1 Hydrogen gas and chlorine gas will react vigorously to produce hydrogen chloride gas. The equation above illustrates this reaction. The reactants, hydrogen and chlorine, are written on the left and the products (hydrogen chloride) on the right. The large number 2 in front of HCl indicates that two molecules of HCl are produced for each 1 molecule of hydrogen and chlorine gas consumed. The 2 in subscript below H indicates that there are two hydrogen atoms in each molecule of hydrogen gas. Finally, the (g) symbols subscript to each species indicates that they are gases. Reacting Species. Species in a chemical reaction is a general term used to mean atoms, molecules or ions. A species can contain more than one chemical element (HCl, for example, contains hydrogen and chlorine). Each species in a chemical equation is written: formula_2 E is the chemical symbol for the element, x is the number of atoms of that element in the species, y is the charge (if it is an ion) and (s) is the physical state. For example, ethyl alcohol would be written formula_3 because each molecule contains 2 carbon, 6 hydrogen and 1 oxygen atom. A magnesium ion would be written formula_4 because it has a double positive ("two plus") charge. Finally, an ammonium ion would be written formula_5 because each molecule contains 1 nitrogen and 4 hydrogen atoms and has a charge of 1+. Coefficients. The numbers in front of each species have a very important meaning—they indicate the relative amounts of the atoms that react. The number in front of each species is called a coefficient. In the above equation, for example, one H2 molecule reacts with one Cl2 molecule to produce two molecules of HCl. This can also be interpreted as "moles" (i.e. 1 mol H2 and 1 mol Cl2 produces 2 mol HCl). It is important that the "Law of Conservation of Mass" is not violated. There must be the same number of each type of atoms on either side of the equation. Coefficients are useful for keeping the same number of atoms on both sides: formula_6 If you count the atoms, there are four hydrogens and two oxygens on each side. The coefficients allow us to balance the equation; without them the equation would have the wrong number of atoms. Balancing equations is the topic of the next chapter. Other Information. Occasionally, other information about a chemical reaction will be supplied in an equation (such as temperature or other reaction conditions). This information is often written to the right of the equation or above the reaction arrow. A simple example would be the melting of ice. formula_7, which could be written as formula_8 Reactions commonly involve catalysts, which are substances that speed up a reaction without being consumed. Catalysts are often written over the arrow. A perfect example of a catalyzed reaction is photosynthesis. Inside plant cells, a substance called "chlorophyll" converts sunlight into food. The reaction is written: formula_9 

Nouns are words that represent something perceived or conceived, like an apple or a thought. In French, all nouns have a grammatical gender; that is, they are either masculine (m) or feminine (f). Most nouns that express people or animals have both a masculine and a feminine form. For example, the two words for "the cat" in French are "le chat" (m) and "la chatte" (f). However, there are some nouns that talk about people or animals whose gender is fixed, regardless of the actual gender of the person or animal. For example, is always feminine, even when it's talking about your uncle; is always masculine, even when it's talking about your female professor or teacher. The nouns that express things without an obvious gender (e.g., objects and abstract concepts) have only one form. This form can be masculine or feminine. For example, can only be feminine; can only be masculine. Exceptions. There are many exceptions to gender rules in French which can only be learned. There are even words that are spelled the same, but have a different meaning when masculine or feminine; for example, means "the book", but means "the pound". Some words that appear to be masculine (such as "la photo", which should be masculine but is not because it is actually short for "la photographie") are in fact feminine, and vice versa. Then there are some that just don't make sense; "la foi" is feminine and means "faith" or "belief", whereas "le foie" is masculine and means "liver". In English, the definite article is always "the". In French, the definite article is changed depending on the noun's: There are three definite articles and an abbreviation. "Le" is used for masculine nouns, "La" is used for feminine nouns, "Les" is used for plural nouns (both masculine or feminine), and "L"' is used when the noun is singular and begins with a vowel or silent "h" (both masculine or feminine). It is similar to English, where "a" changes to "an" before a vowel. Unlike English, the definite article is used to talk about something in a general sense, a general statement or feeling about an idea or thing. Elision. "Elision" refers to the suppression of a final unstressed vowel immediately before another word beginning with a vowel. The definite articles "le" and "la" are shortened to "l’" when they come before a noun that begins with a vowel or silent "h". When pronounced, the vowel sound is dropped. Elision does not occur on an aspired "h": In addition to the definite article, elision will also occur with other words, such as "que", "je", "le", "ce", "ne", and "de". The details on these words will be covered in later sections of the book. In English, the indefinite articles are "a" and "an". "Some" is used as a plural article in English. Again, indefinite articles in French take different forms depending on gender and number. The articles "un" and "une" literally mean "one" in French.  "une" is often (more often than not) pronounced ("ewnuh") in poetry and lyric.  "Des fils" does mean "some sons", but is a homograph: it can also mean "some threads" (when pronounced like ). "Some". Note that "des", like "les", is used in French before plural nouns when no article is used in English. For example, you are looking at photographs in an album. The English statement "I am looking at photographs." cannot be translated to French as "Je regarde photographies." because an article is required to tell which photographs are being looked at. If it is a set of "specific" pictures, the French statement should be On the other hand, if the person is just browsing the album, the French translation is Plurality, pronunciation, and exceptions. The plural of most nouns is formed by adding an "-s". However, the "-s" ending is not pronounced. It is the article that tells the listener whether the noun is singular or plural. Most singular nouns do not end in "-s". The "-s" is added for the plural form of the noun. "Fils" is one exception. Whenever the singular form of a noun ends in "-s", there is no change in the plural form. The final consonant is almost always not pronounced unless followed by an "-e" (or another vowel). is also an exception to this rule. Liaison. Remember that the last consonant of a word is typically not pronounced unless followed by a vowel. When a word ending in a consonant is followed by a word beginning with a vowel sound (or silent "h"), the consonant often becomes pronounced. This is a process called "liaison". When a vowel goes directly after "un", the normally unpronounced "n" sound becomes pronounced. Compare the pronunciation to words without liaison: "Une" is unaffected by liaison. Liaison also occurs with "les" and "des". As with elision, an aspired "h" isn't liaised: "Qu’est-ce que c’est ?". To say "What is it?" or "What is that?" in French, is used. To respond to this question, you say "C’est un(e) ["nom"].", meaning "It is a ["noun"]:" Remember that the indefinite article ("un" or "une") must agree with the noun it modifies: "Il y a".  is used to say "there is" or "there are". "Il y a" expresses the existence of the noun it introduces. The phrase is used for both singular and plural nouns. Unlike in English ("is" → "are"), "il y a" does not change form. The "-s" at the end of the most pluralised nouns tells you that the phrase is "there are" instead of "there is". In spoken French, when both the singular and plural forms almost always sound the same, the article (and perhaps other adjectives modifying the noun) is used to distinguish between singular and plural versions. "A" is the present third person singular form of the verb "to have", and "y" is a pronoun meaning "there". The phrase "il y a", then, literally translates to "he has there". This phrase is used in all French tenses. It is important to remember that verb stays as a form of "have" and not "be". "Voici" and "voilà". Like in English, "il y a…" is not often used to point out an object. To point out an object to the listener, use , meaning "over here is/are" or "right here is/are", and , meaning "over there is/are", or "there you have it". 

 "-uire" verbs are conjugated irregularly. Most verbs form the "passé composé" with "avoir", however there are a small number of verbs that are always conjugated with "être". In a general case, these verbs indicate a change in state or position. List of verbs. The verbs that take être can be easily remembered by the acronym "MRS. DR VANDERTRAMP": Direct objects. These verbs take their conjugated "avoir" when they are immediately followed by a direct object For Example, with the direct object "mes bagages" becomes As another example, with the direct object "mes bagages" becomes As another example, but with "ils" instead of "je", with direct object "leur passeport" becomes Subject-past participle agreement. When conjugating with "être", the past participles of the above verbs must agree with the subject of a sentence in gender and number. Note that there is no agreement if these verbs are conjugated with "avoir". Indirect object pronoun - "to it, to them". The French pronoun "y" is used to replace an object of a prepositional phrase introduced by à. Note that "lui" and "leur", and not "y", are used when the object refers to a person or people. Replacement of places - "there". The French pronoun "y" replaces a prepositional phrase referring to a place that begins with any preposition except "de" (for which "en" is used). Note that "en", and not "y" is used when the preposition of the object is "de". Idioms. These verbs are conjugated irregularly, and normally follow the "-er" conjugation scheme. In past participle form, "-ir" is replaced with "-ert" for these verbs. Formation. A common "-rir" verb is "ouvrir": The noun is derived from "ouvrir", and the adjective is derived from its past participle. "-rir" verb exceptions. 1"Mourir" is the only "-ir" verb that takes "être" as its helping verb in perfect tenses (and therefore agrees with the subject as a past participle in a perfect tense). The word is also used as a noun, meaning "death" or "dead person", or as an adjective, meaning "dead": The derived word means "dying" or "person who is dying". "Acquis" is also a noun, meaning "asset". Examples.  

Europe. In 1815, Europe had united to defeat French Emperor Napoleon. For a century since that time, there had been no major war in Europe. Countries had organized themselves in a complex system of alliances. After Napoleon's defeat, the United Kingdom, France, Prussia, Russia, and Austria met in Vienna. These nations decided that if power in Europe was balanced, then no nation would become so powerful as to pose a threat to the others. The most important of these were the German Confederation. In 1871, after defeating France and Prussia, several small German nations combined into the German Empire. This upset the traditional balance of power. German Chancellor Otto Von Bismarck began to construct a web of alliances to protect German dominance. Germany and the United Kingdom were on good terms, as Germany had not built a navy to go up against British sea power. In 1873, Russia, the Austro-Hungarian Empire, and Germany formed the Three Emperors' League. Nine years later, Austria-Hungary, Italy, and Germany formed the Triple Alliance. In 1887, the Reinsurance Treaty ensured that Russia would not interfere in a war between France and Germany. During this time the British Queen Victoria built alliances in her own way. During years of relative peace, she had her children marry into many of the royal families of Europe, believing that this would solidify relations among the nations. In the first decade of the Twentieth century the Kaiser and the King of England were cousins through Victoria, as were the Tsar and Tsarina of Russia. In 1890, Bismarck was fired by Kaiser Wilhelm II, who then began to undo many of Bismarck's policies. He decided to build up a German navy, antagonizing the United Kingdom. He did not renew German agreements with Russia. In 1894, this led Russia to form a new alliance with Germany's rival France. In 1904, France and the United Kingdom decided to end centuries of bitter enmity by signing the Entente Cordiale. Three years later, those two nations and Russia entered the Triple Entente. Imperial Russia began to build its army, as did Germany and Austria-Hungary. War Breaks Out. Austria-Hungary was a patchwork of several nations ruled by the Habsburg family. Several ethnic groups resented rule by the Habsburgs. In June, 1914, the heir to the throne, Archduke Franz Ferdinand, traveled to Sarajevo in Bosnia and Herzegovina. A Serbian nationalist named Gavrilo Princip, who hated Habsburg rule, assassinated the Archduke and his wife. This assassination triggered the First World War. The Austro-Hungarian government decided to retaliate by crushing Serbian nationalism. They threatened the Serbian government with war. Russia came to the aid of the Serbs. To oppose this alliance, Austria-Hungary called on Germany. Kaiser Wilhelm II said his country would give Austria-Hungary whatever it needed to win the war; in effect, a "blank check." In addition to these open agreements, any of these countries might have had secret agreements with other states. The result was almost all of Europe at war, with the largest battlefield ever seen before. In July, 1914, Austria-Hungary declared war on Serbia. Austria-Hungary, Russia, and Germany began to mobilize their troops. The conflict in Austria-Hungary quickly spread over Europe. In August, Germany declared war on France. The Germans demanded that Belgium allow German troops to pass through the neutral nation. When King Albert of Belgium refused, Germany violated Belgian neutrality and invaded. Belgium appealed to the United Kingdom for aid. The British House of Commons threatened that Great Britain would wage war against Germany unless it withdrew from Belgium. The Germans refused, and the United Kingdom joined the battle. The Central Powers, Germany and Austria-Hungary, were pitted against the Allies, the United Kingdom, Russia, and France. The Early Stages. German troops entered Belgium on August 4. By August 16, they had begun to enter France. The French Army met the Germans near the French border with Belgium. France lost tens of thousands of men in less than a week, causing the French Army to retreat to Paris. The Germans penetrated deep into France, attempting to win a quick victory. On August 5, the United States formally declared their neutrality in the war. They also offered to mediate the growing conflict. In the United States, the opinions were divided. Some felt we should aid England, France, and Belgium because they were depicted as victims of barbarous German aggression and atrocities. Others felt we should avoid taking sides. The Allies won a key battle at Marne, repelling the German offensive. The Germans lost especially due to a disorganized supply line and a weak communications network. The French Army, however, had not completely defeated the Germans. Both sides continually fought each other, to no avail. On the Western Front, Germany and France would continue to fight for more than three years without any decisive victories for either side. Meanwhile, on the Eastern Front, Germany faced Russia. In the third week of August, Russian troops entered the eastern part of Germany. Germany was at a severe disadvantage because it had to fight on two different fronts, splitting its troops. However, despite Germany's disadvantage, no decisive action occurred for three years. The United Kingdom used its powerful Royal Navy in the war against Germany. British ships set up naval blockades. The Germans, however, countered with submarines called U-boats. U-boats sank several ships, but could not, during the early stages of the war, seriously challenge the mighty Royal Navy. The war spread to Asia when Japan declared war on Germany in August, 1914. The Japanese sought control of German colonies in the Pacific. Germany already faced a two-front war, and could not afford to defend its Pacific possessions. In October, 1914, the Ottoman Empire entered, allying itself with the Central Powers. The entry of the Ottoman Empire was disastrous to the Allies. The Ottoman Empire controlled the Dardanelles strait, which provided a route between Russia and the Mediterranean. The Ottoman sultan declared holy war- "jihad"- against the Allies. Muslims in the British Empire and French Empire were thus encouraged to rebel against their Christian rulers. However, the Allies' concerns were premature. Few Muslims accepted the sultan's proclamation. In fact, some Muslims in the Ottoman Empire supported the Allies so that the Ottoman Empire could be broken up, and the nations they ruled could gain independence. The Middle Stages. Between 1914 and 1917, the war was characterized by millions of deaths leading nowhere. Neither side could gain a decisive advantage on either front. In 1915, the Germans began to realize the full potential of Submarines. German Submarines engaged in official unrestricted warfare, engaging and sinking any ship found within the war zone regardless of the flag flown. Germany's justification for this use of force was that there was no certain method to ascertain the ultimate destination of the passengers and cargo carried by the ships in the war zone, and thus they were all taken as attempts at maintaining the anti-German blockade. In May, 1915, Italy broke the Triple Alliance by becoming an Allied Power. In October, Bulgaria joined the Central Powers. Each side had induced their new partners to join by offering territorial concessions. Italy prevented Austria-Hungary from concentrating its efforts on Russia, while Bulgaria prevented Russia from having connections with other Allied Powers. In May, 1916, one of the most significant naval battles in World War I occurred. The Royal Navy faced a German fleet during the Battle of Jutland. The Battle proved that the Allied naval force was still superior to that possessed by the Central Powers. The Germans grew even more dependent on U-boats in naval battle. In August, 1916, Romania joined the Allies. Romania invaded Transylvania, a province of the Austro-Hungarian Empire. But when the Central Powers struck back, they took control of important Romanian wheat fields. In 1917, the liberal-democratic government of Russia that was lead by Aleksander Kerensky was over thrown by Vladimir Lenin. When Lenin took over in Russia one of the things he promised was to change world politics. The terms by which Lenin wanted to changed world politics challenged Woodrow Wilson's vision and Lenin's Bolshevik-style revolutions spreading world wide was something that western leaders did not want. The United States Declares War. Through all of this, America was neutral. It adopted the policy of isolationism because it felt that the increasing colonialism in Europe did not affect North America. There was a strong pacifist strain in American society, as evidenced by such popular songs as "I Didn't Raise My Boy to be a Soldier" and "Don't Take My Darling Boy Away," though at the same time many ethnic groups agitated for involvement. Economic links with the Allies also made neutrality difficult. The British were flooding America with new orders, many of them for arms. The sales were helping America get out of its recession. Although this was good for the economic health of the United States, Germany saw America becoming the Allied arsenal and bank. On May 7, 1915, the German navy sunk the Cunard Line passenger ship R.M.S. Lusitania, operating under the flag of Great Britain. Of the 1,959 passengers 1200 died, including one hundred twenty Americans. The ship's quick explosion was due to a hidden cargo of American weaponry, a fact the US government denied. Woodrow Wilson's Secretary of State William Jennings Bryan, a pacifist, resigned over Wilson's responsibility for the placement of arms and the consequent inevitability of war. Various American citizens from different ethnic groups put pressure on their government to join the war. However, the US government was calmed by the Germans, who agreed to limit submarine warfare. In 1917, the Germans reinstated unrestricted submarine warfare in order to cripple the British economy by destroying merchant ships, and break the sea blockade of Britain. On February 24, 1917, the American ambassador received a telegram in London from the British. It enclosed a British-decoded message, originally sent as a ciphered telegram from the German foreign Secretary, Arthur Zimmerman, to the ambassador in Mexico. Zimmerman proposed that the event of the war with the United States, Germany and Mexico would join in alliance. Germany would fund Mexico's conflict with the US: victory achieved. Mexico would then be able to gain their lost territories with Arizona. The message was published in American newspapers on March 1st. On the evening of April 4, 1917 at 8:30 P.M., President Wilson appeared before a joint session of Congress, asking for the declaration of war to make the world "safe for democracy." He was hoping for a quick resolution of the conflict. Congress complied on April 6, 1917. The United States was at last at war with Germany. The last American war had been the minor Spanish-American War a generation before. A draft of men above the age of eighteen followed the call to war, but many more volunteered. Men wanted to escape their lives and join the military for a job and an adventure. The US had to mobilize its military before it could aid the Allies by sending troops. The cadre of the U.S. Army had experience in mobilizing and moving troops from its Mexican expedition, but the Army needed to expand to over one million men, most of which were untrained. This new draft was for a "scientific" army. Military officers were selected by the new "IQ" tests. The new conscripts met "en masse"in embarkation camps in places such as Yahank, New York. Companies exercised together and drilled together. At this time, soldiers were taught techniques such as assembling and disassembling weapons blindfolded. They drilled constantly in formation. There was much amusement at these shaven-headed, first-time solders. The chorus of a contemporary song went, "Good morning, Mister Zip-Zip-Zip,/ With your hair cut just as short as mine,/ Good morning, Mister Zip-Zip-Zip,/ You're surely looking fine!/ Ashes to ashes, and dust to dust,/ If the Camels don't get you,/ The Fatimas must . . . [Camels and Fatimas were brands of cigarettes.] Then they went into ships to be transported to Europe. On the battlefield, they were placed in companies with other Americans, rather than being embedded with other nations. American business and industry became involved as men created more military supplies, and jobs opened for building and designing new materials to be used in battle. However, for purposes of the battle, most of the weapons American soldiers carried had been designed by Europeans. An American-style tank was produced but only came out after the war. In the same way, the Navy could send a battleship division to assist the British Grand Fleet, but needed to expand. To supply the American forces, new supply lines in France would be needed south of the British and French lines, which meant the U.S. would take over the southern part of the Western Front battle line. The US could and did help the Allies with monetary assistance. Increased taxes and the sale of war bonds allowed the US to raise enormous sums of money. Politicians and celebrities, as well as such movie stars as Charlie Chaplin and Mary Pickford, headed huge patriotic "Bond Rallies," where people were encouraged to buy bonds. A government committee to influence the public on the war was formed, the Committee on Public Information or CPI. Among its organs of publicity were the "Four Minute Men," speakers who talked on pertinent subjects on Vaudeville stages, in movie theaters, and in public assemblies. There was also an organization of private citizens formed to root out German sympathizers, the American Protective League. A hit movie, "The Kaiser, the Beast of Berlin," was one hit of 1918. To strengthen the United States in this time of stress, families were encouraged to grow Victory Gardens, and American women and African Americans were encouraged to go into jobs the servicemen had left. This is the beginning of the Great Migration, when Southern African Americans began moving to Northern cities for jobs. Trench Warfare. The U.S. commander, General John J. "Black Jack" Pershing, faced immense pressure from the British and French governments to use American forces in small units to reinforce depleted British and French units. This was impossible politically. Pershing insisted to General Foch, the Generalissimo of the Allied armies, that the U.S. Army would fight as a single Army. Pershing did not want to give his men to other Allied commanders, many of whose strategies he disagreed with. The European method of fighting, as it had been since the Boer War, was trench warfare. An army on the French battleground protected itself from the enemy with zigzag trenches, mines, barbed wire, and a line of rifles and machine guns. Between the enemy lines was a contested area, "no man's land." An attempt to advance toward the enemy was met with gunfire. These trenches stalemated military advances, as any man who raised his head from the trenches would be shot. At the Battle of the Somme in 1916, for example, Allied troops suffered 600,000 dead and wounded to earn only 125 square miles; the Germans lost 400,000 men. Rain fell in the trenches. As one song put it, a soldier was "Up to your waist in water, up to your eyes in slush." The damp produced a foot disease known as "trench foot": if untreated, it could rot flesh from the bone. The close, unsanitary conditions of the front lines encouraged fleas and lice, and typhoid, typhus, and dysentery caused deaths unrelated to gunfire. Worst of all, perhaps, was that sometimes the enemy would gas your trench. There was no way to escape, and sometimes no masks to protect you. First used by the Germans in April 1915, chlorine gas stimulated overproduction of fluid in the lungs, leading to death by drowning. One British officer tended to troops who had been gassed reported that, Chlorine, mustard, and phosgene gas would continue in use throughout the war, sometimes blistering, sometimes incapacitating, and often killing. The End of the War. Despite the fact that the Germans could concentrate their efforts in one area, the Central Powers faced grim prospects in 1918. Encouraged by the United States joining the war, several nations joined the Allied Powers. The four Central Powers of Germany, the Austro-Hungarian Empire, the Ottoman Empire, and Bulgaria faced the combined might of the Allied Powers of the United Kingdom and the British Empire, Australia, New Zealand, Canada, South Africa, France, Belgium, Japan, Serbia, Montenegro, San Marino, Italy, Portugal, Romania, the United States, Cuba, Panama, Guatemala, Nicaragua, Honduras, Haiti, Costa Rica, Brazil, Liberia, Siam (Thailand) and China (some of the above nations did not support the war with troops, but did contribute monetarily.) The Germans launched a final, desperate attack on France, but it failed miserably. Due to Allied counterattacks, the Central Powers slowly began to capitulate. Bulgaria was the first to collapse. A combined force of Italians, Serbs, Greeks, Britons, and Frenchmen attacked Bulgaria through Albania in September, 1918. By the end of September, Bulgaria surrendered, withdrawing its troops from Serbia and Greece, and even allowing the Allies to use Bulgaria in military operations. British forces, led by T. E. Lawrence (Lawrence of Arabia), together with nationalist Arabs, were successful in the Ottoman Empire. About a month after Bulgaria's surrender, the Ottoman Empire surrendered, permitting Allies to use the Ottoman territory, including the Dardanelles Strait, in military operations. The Austro-Hungarian Empire also decided to surrender in October. The royal family, the Habsburgs, and the Austro-Hungarian government desperately sought to keep the Empire of diverse nationalities united. Though Austria-Hungary surrendered, it failed to unite its peoples. The once-powerful Austro-Hungarian Empire was destroyed by the end of October, splitting into Austria, Hungary, Czechoslovakia, and Yugoslavia. Germany, remaining all alone, also decided to surrender. President Wilson required that Germany accede to the terms of the Fourteen Points, which, among other things, required Germany to return territory acquired by the Treaty of Brest-Litovsk to Russia and the provinces of Alsace and Lorraine to France. Germany found the terms too harsh, while the Allies found them too lenient. But when German Emperor Wilhelm II abdicated the throne, the new German government quickly agreed to Wilson's demands. On November 11, 1918, World War I had come to an end. The war had been marked by millions and millions of casualties. The deaths were so wide-spread and so vast that people in England talked of "The Lost Generation." Many died in battle, others died from disease and some even died after when hit with an influenza that spread throughout the whole world in 1918. Destruction of factories and farms, not to mention houses, created economic damage, and was one of the factors creating widespread European starvation during the winter of 1918-1919. In contrast, damage was low in America. Although America had also suffered from the flu, the war had not touched U.S. shores. For years afterward, in Germany, France, and even British Commonwealth countries such as Canada, you could see men wearing artificial tin faces to hide battle wounds, men who wheezed because of damage from poison gas, and "war cripples" begging with bowls on the street corner. But American men were by-and-large intact. U.S. factories had been fully supplied, and the nation was on a sounder financial footing because of profit from the War. At the end of the war, as American soldiers returned from Europe, employment rose. Some of the veterans returned to find themselves without homes or jobs. Overseas, some Black men were organized into a top fighting unit. However, Homeland outrage at the success helped to fuel riots in the notorious "Red Summer" of 1919. Each veteran returned with a certificate promising certain monies for their service; however, the certificate could not be cashed in until 1945. American After-Effects Of The War. Suffrage For American Women. An after-effect of the employment of women during the war was the Nineteenth Amendment, giving them the right to vote. There had been a Suffrage movement since the nineteenth century. President Wilson and his contemporaries were reluctant, but conceded it as a "quid pro quo". After the servicemen had returned, there were still two million more women in the workforce. However, instead of being in factory jobs, they were largely only permitted in "women's work": the "caring professions" such as nurses or teachers, secretaries or "stenographers," and waitresses, cooks, or washerwomen. These jobs paid little, and women were often expected to quit the job when they got married. Politically active women still remained excluded from local and national power structures. Their voluntary organizations used tactics that advanced modern pressure-group politics. Issues ranged from birth control, peace, education, Indian Affairs, or opposition to lynching. Women in these associations lobbied legislators to support their causes. At the state level women achieved rights such as the ability to serve on juries. Increasing Racial Tension. A force of African Americans had served in the War: Fighting units such as the Harlem Hellfighters proved their mettle. Yet they received no recognition. Upon their return, particularly to the South, they faced resentment and sometimes lynchings. In 1913, President Woodrow Wilson had segregated the Civil Service, which before had been an employer for black Americans. The war and the Great Migration triggered more oppression and more violence. When Caucasian White Americans were drawn into the army or defense industries, their jobs were sometimes given to black workers, for lower wages. The owners considered this a double good, keeping production going while destroying the workingman unions which agitated for higher pay. But while bosses were dependent, for the moment, on black workers, they did not think they owed anything to these employees. An influx of unskilled Black strikebreakers into East St Louis, Illinois, heightened racial tensions in 1917. Rumors that Blacks were arming themselves for an attack on Whites resulted in numerous attacks by White mobs on Black neighborhoods. On July 1, Blacks fired back at a car whose occupants they believed had shot into their homes and mistakenly killed two policemen riding in a car. The next day, a full scale riot erupted which ended only after nine Whites and 39 Blacks had been killed and over three hundred buildings were destroyed. The anxieties of war helped fuel violence: Southern and Midwestern war vigilance committees formed the matrix for a revival of the Ku Klux Klan. The Great Experiment. On August 1, 1917, the Senate voted to send the Eighteenth Amendment to the Constitution to the states for ratification. The vote was bi-partisan, 65 to 20. Section One read, in part, By 1919, the requisite number of states had ratified the Amendment. The Amendment actually came into effect, under its own terms, one year after ratification. The Temperance movement had been in effect for more than a hundred years, since the Second Great Revival. The late 19th and early 20th century Anti-Saloon League had been successful in turning the discussion from social discouragement of alcohol to legal prohibition of the substance. It was not simply a creature of Protestant or Catholic Churches, but was "united with Democrats and Republicans, Progressives, Populists, and suffragists, the Ku Klux Klan and the NAACP, the International Workers of the World, and many of America's most powerful industrialists including Henry Ford, John D. Rockefeller, Jr., and Andrew Carnegie – all of whom lent support to the ASL's increasingly effective campaign." But if Prohibition was not a simple reaction to World War I, it drew strength from the conflict. The War was funded by an income tax, thus unlocking saloon profits from the interests of the nation. The Lever Act of 1917, with the aim of feeding soldiers, prohibited grain from being used for alcoholic beverages. There were nationalist concerns, too: the most famous brewers of beer had German last names. Treaty of Versailles. The Treaty of Versailles was the peace settlement signed after World War One had ended in 1918 and in the shadow of the Russian Revolution and other events in Russia. The treaty was signed at the vast Versailles Palace near Paris - hence its title - between Germany and the Allies. The three most important politicians there were David Lloyd George, Georges Clemenceau, and Woodrow Wilson. The Versailles Palace was considered the most appropriate venue simply because of its size - many hundreds of people were involved in the process and the final signing ceremony in the Hall of Mirrors could accommodate hundreds of dignitaries. Many wanted Germany, now led by Friedrich Ebert, smashed. Others, like Lloyd George, were privately more cautious. On June 28th 1919, the chief Allied and Associated Powers of the United Kingdom, the United States, France, Italy, and Japan met with the Central Powers in France to discuss a peace settlement. British Prime Minister David Lloyd George, American President Woodrow Wilson, and French President Clemenceau were known as "The Big Three." Each of the Allied and Associated Powers had distinct national interests. The UK wanted to keep the Royal Navy supreme by dismantling the German Navy. The British also wished to end Germany's colonial empire, which might have become a threat to the vast British Empire. Lloyd George wanted to be hard on the Germans to bolster his popular support at home in Great Britain.. Italy wanted the Allies to fulfill their promise of territory given to Rome at the beginning of the war. Clemenceau wanted Germany to be brought to its knees so it could never start another war against France. The French also wanted Germany to compensate Paris for damage inflicted on France during the War. Japan had already largely served its interests by taking over German colonies in the Pacific. President Wilson's main goal for the conference was the creation of a "League of Nations." He believed that such an organization was essential to preventing future wars. Many historians believe that Wilson's concentration on the League, forcing him to sacrifice possible compassion toward Germany, helped contribute to the conditions leading to World War II. The Treaty of Versailles forced Germany to cede Alsace and Lorraine to France, dismantle its Army and Navy, give up its colonial Empire, pay massive reparations to the Allies, and take full responsibility for causing the war. The conference also led to the creation of the League of Nations. The US Senate, however, did not consent to the Treaty, and the European powers were left to enforce its provisions themselves. This eventually led to violations of the treaty by Germany, which then led to the Second World War. The treaty crippled the Weimar Government in Berlin and led to great bitterness in Germany, which helped to strengthen Adolf Hitler's National Socialist, or Nazi Party. Questions For Review. 1. What extended a conflict between Serbians and the Austro-Hungarian Empire into the globe-straddling World War I? 2. How did the advances of technology lead to trench warfare? 3. What in 1914-1915 led to an economic advantage for America from the War? What were the nation's post-war advantages? 4. Listen to songs from this period: "I Didn't Raise My Boy to be a Soldier." "Oh, How I Hate to Get up in the Morning." "Oh! It's a Lovely War." "Over There." "K-k-k-Katy." What is the song's point of view? 



Introduction. In this section we look at the polynomial in some commutative ring with identity. What is interesting is that studying polynomials over some commutative ring with identity acts very much like numbers; the same rules often are obeyed by both. Definitions. A "polynomial" over some commutative ring with identity R is an expression in the form and n ∈ N, and "x" is some indeterminate ("not" a variable). Terminology. Given the first nonzero term in the polynomial, i.e. the term "a"nxn above: In the above, if "a"i=0 for all i, the polynomial is the "zero polynomial". Properties. Let R[x] be the set of all polynomials of all degrees. Clearly R is closed under addition and multiplication (although in a non-straightforward way), and thus we have that R[x] is itself a commutative ring with identity. Assume now R is a field F; we do this so we can define some useful actions on polynomials Division algorithm. Firstly recall the division algorithm for numbers, that each number can be decomposed into the form where "q" is the quotient and "r" the remainder and r&lt;n. Now, since we have that F is a field, we can do something similar with the polynomials over F, F[x]. If f("x"), g("x") ∈ F[x], with g("x") nonzero: Again, q("x") is known as the quotient polynomial and r("x") the remainder polynomial. Furthermore, we have the degree of r(x) ≤ degree of f(x) We perform divisions by polynomial long division. For brevity we omit the "x"k terms. Here's an example. We divide "x"3+"x"+2 by "x"-1. First write: Note we place a 0 in any polynomial not present. Now "x"3/"x" = "x"2, so we place a 1 in the second column to get Multiply "x"2 throughout the divisor "x"-1 to get "x"3-"x"2, which is (1 -1), so write this below like the following: Now subtract (1 0) and (1 -1), drop the third 1 to get: Now repeat, but divide by "x"2 now (since we have subtracted and gotten (1 1) - "x"2 + "x"), and continue in the same fashion, to get: So the quotient is "x"2+"x"+2, and the remainder is 4. Euclidean algorithm. Now we have a working division algorithm for polynomials, the Euclidean algorithm, and hence the gcd of two polynomials can readily be found. Examples. Let's use a similar example above: what is gcd("x"3+"x"+2, "x"-1)? We've shown already above that "x"3+"x"+2=("x"2+"x"+2)("x"-1) + 4 Proceeding in the normal fashion in the Euclidean algorithm and the greatest common divisor of any monomial and an integer is clearly 1, so "x"3+"x"+2 and "x"-1 are coprime. For a second example, consider gcd("x"2-1,"x"2+2"x"+1) Since factors of -1 make no difference, gcd("x"2-1,"x"2+2"x"+1) is -("x"+1) Irreducibles. We've seen that "x"3+"x"+2 and "x"-1 are coprime; they have no factors in common. So, are we able to determine "prime" polynomials? Indeed we can - depending on the field that the polynomial lies in. We call these "irreducibles" instead of primes. Example. Take p("x")="x"3 + "x"2 + 2 over Z3. Now we can factor this polynomial if it has a root - from the factor theorem (which also holds for polynomials over any commutative ring with identity) p(k)=0 means k is a root. So, let's look at the following: Since we're in Z3, we luckily only need to check three values p(0)=2 p(1)=1 p(2)=2 So we have p("x") having no roots - it is irreducible ("prime"). Now observe an interesting fact. Take the exact same polynomial but instead over Z2. The polynomial then is equivalent to and thus has root p(0)=0 and thus "is reducible" but "over" Z2 So the reducibility of the polynomial depends on the field it is in. Showing irreducibility. The general procedure to show a polynomial is irreducible is: effectively a proof by cases. For example, consider the polynomial q("x")="x"4+"x"+2 in Z3. To prove it is irreducible, observe that q("x") could be factorized in the following ways: So we can prove 1, 2, 3 by showing it has no linear factors. 4 is a little more difficult. Let us proceed to show it has no linear factors: Observe So q has no linear factors. Now, we need to show that q is not the product of two irreducible quadratics. In Z3, we have the quadratics We can identify the irreducible quadratics easily by inspection. We then obtain If we can show that neither of these polynomials divide q("x")="x"4+"x"+2, we have shown q("x") is irreducible. Let us try "x"2+1 first. We have a remainder, so "x"2+1 doesn't divide q. On dividing the other polynomials, we all get a remainder. ("Verify this for yourself as practice"). So q("x") is irreducible in Z3. Modular arithmetic and polynomials. Since we have a working polynomial division and factor theorem, and that polynomials appear to mimic the behaviour of the integers - can we reasonably define some sort of modular arithmetic with polynomials? We can indeed. If we have a field Zp["x"] and we wish to find all the remainders (remember, these remainders are polynomials) on dividing by some polynomial m("x"), we can do so by polynomial long division. If m("x") is irreducible, then the set of remainders as above forms a field. 

Category theory is the study of "categories", which are collections of objects and "morphisms" (or arrows), from one object to another. It generalizes many common notions in Algebra, such as different kinds of products, the notion of kernel, etc. See Category Theory for additional information. Definitions &amp; Notations. Definition 1: A "(locally small) category" formula_1 consists of These obey the following axioms: Note that we demand neither formula_2 nor formula_3 to be sets; if they are both in fact sets, then we call our category "small". Definition 2: A morphism formula_13 has associated with it two functions formula_27 and formula_28 called "domain" and "codomain" respectively, such that formula_11 if and only if formula_30 and formula_31. Thus two morphisms formula_32 are composable if and only if formula_33. Remark 3: Unless confusion is possible, we will usually not specify which Hom-set a given morphism belongs to. Also, unless several categories are in play, we will usually not write formula_34, but just "formula_7 is an object". We may write formula_36 or formula_37 to implicitly indicate the Hom-set formula_13 belongs to. We may also omit the composition symbol, writing simply formula_39 for formula_40. Basic Properties. Lemma 4: Let formula_7 be an object of a category. The identity morphism for formula_7 is unique. "Proof": Assume formula_43 and formula_44 are identity morphisms for formula_7. Then formula_46. Example 5: We present some of the simplest categories: Initial and Final Objects. Definition An object formula_63 in a category is called initial or cofinal, if for any object formula_7 there exists a unique morphism formula_65 Lemma If formula_63 and formula_67 are initial objects, then they are isomorphic. "Proof": Let formula_68 and formula_69 be the unique morphisms between formula_70 and formula_67. Given that both formula_70 and formula_67 have a unique endomorphism because of their initiality, this morphism must be the identity. Therefore formula_74 and formula_75 are the respective identity morphisms, making formula_70 and formula_67 isomorphic. Definition An object formula_78 in a category is called final or coinitial, if for any object formula_7 there exists a unique morphism formula_80 Lemma If formula_81 and formula_82 are final objects, then they are isomorphic. "Proof" Pass to isomorphicness of initial objects in the cocategory. Some examples of categories. In all the examples given thus far, the objects have been sets with the morphisms given by set maps between them. This is not always the case. There are some categories where this is not possible, and others where the category doesn't naturally appear in this way. For example: 

President Grover Cleveland. In 1884, as the presidential campaign season approached, the Republican party chose former Speaker of the House James G. Blane as its candidate, with John Logan as the vice presidential candidate. Against them the Democrats ran New York governor Stephen Grover Cleveland for a presidential candidate and vice presidential candidate Thomas A. Hendricks. Cleveland and Hendricks won with the combined support of Democrats and reform Republicans, the "Mugwumps." Grover Cleveland was the only Democrat elected to the presidency during the era of Republican political domination that lasted from 1860 to 1912. In fact, he won the popular vote for president three times, in 1884, 1888, and 1892. His last election in 1892 defeated Republican president Benjamin Harrison. He was thus the only president to serve two non-consecutive terms, and is the only individual to be counted twice in the numbering of the presidents. (In that last election, Cleveland's vice president was Adlai E. Stevenson; Harrison's Vice President was Whitelaw Reid.) Cleveland's conservative economic stand in favor of the gold standard brought him the support of various business interests. The democrats then won control of both houses of Congress. Racism. Racism was a major blight of the 1890s', largely unacknowledged by the government in Washington. The freedoms which had been given to the former slaves by Emancipation were taken away in Southern states. Every one of them passed laws which disenfranchised African Americans, including poll taxes and literacy tests. The notorious Jim Crow laws, named for an old minstrel show song, started to take effect even in some states which had not seceded. Blacks were barred from public drinking fountains, bathrooms, and train cars, and directed by sign to "Negro only" facilities which were often dirty and defective. A new set of photographic post cards started going through the mail. These showed the results of public lynchings, largely of African American men. Many of these postcards show large crowds, some including White children, and the hung, mutilated, or burnt body of the victim. These proud shows seldom resulted in any prosecution of the attackers, who sometimes included local public officials. In 1898 white citizens of Wilmington, North Carolina, resenting African Americans’ involvement in local government and incensed by an editorial in an African American newspaper accusing white women of loose sexual behavior, rioted and killed dozens of blacks. In the fury’s wake, white supremacists overthrew the city government, expelling black and white office holders, and instituted restrictions to prevent blacks from voting. Industrialization. Rise of Industrial Power. In the 1870's, as the Civil War receded into memory, the United States became a leading Industrial power. Advances in technology and new access to the immense resources of the North American continent drove American Industrialization. This industrialization brought the growth of new American cities such as Chicago, and the arrival of a flood of immigrants from all over Europe to man the factories. During the Gilded Age, businessmen reaped enormous profits from this new economy. Powerful tycoons formed giant trusts to monopolize the production of goods that were in high demand. Andrew Carnegie built a giant steel empire using vertical integration, a business tactic that increased profits by eliminating middlemen from the production line. Jay Gould grabbed the railroads, and then the resources brought by those railroads. Though industrialization caused many long-term positives, it did cause problems in the short-term. Rich farmers who could afford new machinery grew even richer, while poorer farmers were forced to move to urban areas, unable to compete in the agricultural sector. In 1878 the U.S. had entered an era of success after a long downfall of the mid 1870's. The number of manufacturing plants and number of people doubled. By the 1900s the South had more than 400 mills. Women and children worked in bad conditions for up to 12 to 16 hours per day. They only made about a half a dollar per day, equivalent to fifteen 2015 dollars. The War of the Currents. During the 1880's and 1890's the War of the currents was fought between American Thomas Alva Edison and George Westinghouse, then Titans in the American electrical industry. Type Revolution. In 1868 the typewriter was perfected by an editor named Christopher Sholes. This invention brought about a wave of new employment opportunities for women. The machine was made popular by several authors, most notably Mark Twain: the first to make a typewritten manuscript, and the first to send it to a printer. Not even the best and most accurate copperplate printer could fit as many words on a page as the standard typewriter. (For many years, a "typewriter" was the name for the person operating the machine.) It was cheaper to employ women than men as typists and telephone and telegraph operators. Along with this new machine also came other inventions such as the telephone and the telegraph. In the 1890s the number of telephone and telegraph operators went up 167 percent, and the number of women stenographers and typists went up 305 percent. Internal Combustion Engine. Among the early innovations in technology was development of the internal-combustion engine. In 1885 a German engineer, Gottlieb Daimler, built a lightweight engine driven by vaporized gasoline. This development inspired American Henry Ford. In the 1880's, while he was still an electrical engineer in Detroit's Edison Company, Ford experimented in his spare time using Daimler's engine to power a vehicle. George Selden, a Rochester, New York, lawyer, had already been tinkering with such technology, but it was Ford who created a massive industry. Factory Jobs. As industrialization increased, more job opportunities opened. Factory jobs were perfect for women and children, with their smaller hands and their lower pay rates. Despite terrible work conditions, increasing numbers of women moved from the home to factories. But while women became part of the factory floor, virtually none were trusted with management, or even with handling money. The factories also took in immigrants and used them as cheap labor. Catholic immigrants from Ireland and Germany and Jews from Eastern Europe were second-class citizens in the workplace, with very low wages and no benefits. Without safety precautions, workers often suffered serious injury and lost their jobs. Workers adjusted to mechanization as best they could. Some people submitted to the demands of the factory, machine, and time clock. Some tried to blend old ways of working into the new system. Others turned to resistance. Individuals challenged the system by ignoring management's orders, skipping work, or quitting. But also, anxiety over the loss of independence and a desire for better wages, hours, and working conditions drew disgruntled workers into unions. In the cities, laborers and employers often clashed over wages, sanitary conditions, working hours, benefits, and several other issues. Laborers organized themselves into unions to negotiate with companies. The companies, however, attempted to shut down labor unions. Some imposed "yellow dog contracts", under which an employer could dismiss a worker who participated in union activity. In 1886, the American Federation of Labor was formed to fight for laborers in general. The AFL and other union groups employed as many tactics as possible to force employers to accede to their demands. One tactic was the strike. Some strikes escalated into riots, as with the Knights of Labor's strike in 1886 becoming the Haymarket Riots. The Haymarket Riots of 1886 occurred when an unknown person threw a dynamite bomb into a group of police officers. Eight officers were killed in the explosion and gunfight that ensued. As a result, eight anarchists were tried for murder -- four were sentenced to death and one committed suicide. Significant Strikes. The Pullman Strike occurred in 1894, in response to Pullman Company workers' wages being cut following the Panic of 1893, an economic depression which was caused in part by excessive railroad speculation. Approximately 3,000 workers began the strike on May 11. Many of the workers were members of the American Railway Union, and although the strike began without authorization from union officials (known as a "wildcat strike"), the ARU eventually supported the strike by launching a nationwide boycott of Pullman cars on June 26. Within four days, approximately 125,000 ARU members had quit their jobs rather than switch Pullman cars. On July 6, President Cleveland sent Army troops to break up the strike, ostensibly because it prevented delivery of mail and was considered a threat to public safety. The companies sometimes retaliated against strikes by suing the unions. Congress had passed the Sherman Antitrust Act to prevent trusts, or corporations that held stock in several different companies, from obstructing the activities of competitors. Though the Sherman Act was intended to target trusts, the companies sued the union under it, claiming that unions obstructed interstate commerce. During the machine age, there were a number of strikes that took place due to the demands from factories and time clocks. It was hard for individuals to adjust to that system, and as a result, they challenged the system by ignoring management's orders, skipping work, or quitting. The desire and longing for better wages let to anxiety and frustration. Like farming and mining, industry was massive in size and changed not only the nature of the work but the person doing it. Soon, all of these disgruntled individuals formed specialized groups into unions. The different jobs varied in not only skill, but other things as well that were non-related to worker conflict; race, sex, etc. These jobs were such as working on/in railroads, steel factories, and automobiles. The outcome for many working in labor during the Gilded Age led to horrific labor violence. Industrialists and workers literally fought over control of the workplace. Many suffered due to the strikes and riots and it inevitably led to deaths, loss of jobs, and often continuous violence. For most American workers, the Machine age had varying results. At times there was no job stability and when costs of living would increase drastically there were even more problems. Prices and Wages Fall. Prices, and consequently wages, fell sharply in about the 1870s and stayed that way all the way through the 1970s. The prices of necessities in the late 1800s were: 4 pounds butter for $1.60, 1 bag of flour $1.80, a quart of milk for $0.56, vegetables $0.50, 2 bushels of coal $1.36, soap, starch, pepper, salt, vinegar, etc. $1.00, rent for $4.00 a week, and more. The average total of a person's wages was $16.00. By the time that person bought the necessities such as food and soap and rent, most, if not all, of the money would be gone. Woman's Movement. The Woman's Movement, the group of women advocating women's suffrage and equality, continued after the Civil War. Even many women who were not interested in the Vote were creating clubs and crusades, advocating for public issues both before and after marriage. The argument of "separate spheres" for men and women, the rough public world for the former and the gentle domestic world for the latter, was being contested. Jane Addams argued that “If women have in any sense been responsible for the gentler side of life which softens and blurs some of its harsher conditions, may not they have a duty to perform in our American cities?” Urbanization. With industrialization came urbanization. The increasing factory businesses created many more job opportunities in the cities. Soon people began to flock from rural, farm areas, to large cities. Minorities and immigrants added to these numbers. Factory jobs were the only jobs some immigrants could get, and as more came to the cities to work, the larger the urbanization process became. In 1870 there were only two American cities with a population of more than 500,000, but by 1900 there were six, and three of these, New York, Chicago, and Philadelphia had over one million inhabitants. Roughly 40 percent of Americans lived in cities and the number was climbing. These large populations in the cities caused the crime rates to go up, and disease was rapidly spreading. Not only did urbanization cause cities to grow in population, it also caused cities to grow in building size. Skyscrapers were being built in the cities and the idea of mass transit had started. With these mass transits being built it allowed people to commute to work from further distances. Suburbs were beginning to form and higher class families began to move to them to get out of the over crowded city but still gave them the ability to go into the city to work each day. City living was for the lower class the upper class had enough money to get away from all of the pollution and the city stench. This still holds true today in larger cities a lot of the nicer homes are located further out from the center of the city. For example, in the city of Chicago, you will find a lot of the nicer homes away from the city, and more towards the suburbs. In this case, this is because there are a lot of violence in the inner city. Therefore, people try to live more further out from the city in order to stay away from the violence. Agriculture. In the late 1880s and early 1900s, a typical farm would be about 100 acres. Farmers could only plow with the aid of a horse or a mule. Later on the internal combustion engine was used to create tractors. Unlike Southern cotton plantations, most farms raised a variety of foodstuffs, breeding cows, pigs and chickens, and growing turnips, potatoes, carrots, wheat, and corn. They were often self-sufficient. Farmers made their bread from their own wheat, and killed the runt pig for their own table. While industry generally increased in importance, farmers struggled due to debt and falling prices. In the 1880s there were crop failures. Steamships and railways brought in wheat from abroad, lowering American farm prices still more. Economic transformation created industrial prosperity and new lifestyles, but in states still dominated by farming these changes also had a widespread negative effect. Crop diversification and a greater focus on cotton as a cash crop did not give many farmers any potential to get ahead. American farmers helped to create regulation of the railroads. When domestic farmers needed to transport their crops, they also had to rely on the railroad system. But railroads often charged outrageous prices. Farmers, small merchants, and reform politicians started to demand rate regulation. In 1877, in Munn v. Illinois, the Supreme Court upheld the principle of state regulation, declaring that grain warehouses owned by railroads acted in the public interest and therefore must submit to regulation for the common good. By 1880 fourteen states had established commissions to limit freight and storage charges of state-chartered lines. Between 1860 and 1905 the number of farms tripled from two million to six million. In 1905 the number of people living on farms grew to thirty-one million. The value of farms went from eight billion in 1860 to thirty billion in 1906. Then as now, wheat was a major crop, creating such common food as bread, a major source of both starch and protein for poorer people. Farmers had to rise early, often at four or five in the morning. Cows and goats had to be milked twice a day, at morning and at evening. Chickens' eggs were gathered every morning, cleaned, and packed in cases. Because they laid eggs, female chickens or pullets were more important than the male chickens or roosters. Because of this, and to keep roosters from attacking each other, poultry farmers would have only one rooster with several hens. After the Civil War, more prosperous farmers gained more machinery to plant and harvest their crops. In 1879 the centrifugal cream separator was patented. In 1885, chicken raising became a lot more profitable due to the invention of the mechanized incubator. Complicated horse-drawn mechanical combines and threshers were used about this time. With the aid of machines a farmer could harvest about 135 acres of wheat; without them, he or she could only harvest about 7.5 acres of wheat in the same amount of time. The Montgomery Ward &amp; Company mail order catalogue of 1895 listed various grist mills, seeders and planters, a hay pitcher, a hay tedder, and nearly a full page of mechanical churns. Imperialism. As time progressed, Industrialization caused American businessmen to seek new international markets in which to sell their goods. In addition, the increasing influence of Social Darwinism led to the belief that the United States had the inherent responsibility to bring concepts like industry, democracy and Christianity to less scientifically developed, "savage" societies. The combination of these attitudes, along with other factors, led the United States toward Imperialism, the practice of a nation increasing its sphere of influence. The Orient. In the Orient, Russia, Japan, the United Kingdom, France, and Germany all exercised influence. US Secretary of State John Hay endorsed the "Open Door Policy", under which all foreign powers would exercise equal economic power in the Orient. The US thus protected its interests in China and maintained a balance of power there. Chinese nationalists known as the "Righteous Fists of Harmony", or "Boxers" in English, who resented foreign influence, promoted hatred of non-Chinese as well as Chinese Christians. In June 1900 in Beijing, Boxer fighters threatened foreigners and forced them to seek refuge in the Legation Quarter. In response, the initially hesitant Empress Dowager Cixi, urged by the conservatives of the Imperial Court, supported the Boxers and declared war on foreign powers. Diplomats, foreign civilians, soldiers, and Chinese Christians in the Legation Quarter were under siege by the Imperial Army of China and the Boxers for 55 days. The siege was raised when the Eight-Nation Alliance brought 20,000 armed troops to China, defeated the Imperial Army, and captured Beijing. The Boxer Protocol of 7 September 1901 specified an indemnity of 67 million pounds (450 million taels of silver), more than the government's annual tax revenue, to be paid over a course of thirty-nine years to the eight nations involved. Spanish Territories. By 1825 Spain had acknowledged the independence of its possessions in the present-day United States. The only remnants of the Spanish Empire in the Western Hemisphere were Cuba, Puerto Rico, across the Pacific in the Philippines Islands, as well as the Carolina, Marshall, and Mariana Islands (including Guam) in Micronesia. In 1898, the American battleship USS "Maine" was destroyed by an explosion in the Cuban Harbor of Havana. Although later investigations proved that an internal problem was to blame, at the time it was thought that Spanish forces had sunk it. On the advice of Assistant Secretary of the Navy Theodore Roosevelt, President William McKinley asked Congress to declare war on April 11, 1898. Senator Henry M. Teller of Colorado added an amendment to the proposed U.S. declaration of war against Spain on April 19, which proclaimed that the United States would not establish permanent control over Cuba. The amendment stated that the United States "hereby disclaims any disposition of intention to exercise sovereignty, jurisdiction, or control over said island except for pacification thereof, and asserts its determination, when that is accomplished, to leave the government and control of the island to its people." At that time Spanish troops stationed on the island included 150,000 regulars and 40,000 irregulars and volunteers while rebels inside Cuba numbered as many as 50,000. Total U.S. army strength at the time totalled 26,000, requiring the passage of the Mobilization Act of April 22 that allowed for an army of at first 125,000 volunteers (later increased to 200,000) and a regular army of 65,000. On April 25, 1898 Congress declared war on Spain. The United States Navy won two decisive naval battles, destroying the Spanish Pacific Fleet at Manila in the Philippines and the Atlantic fleet at Santiago, Cuba. The U.S. then landed forces in Cuba, which fought the tropical climate and associated diseases as well as the Spanish forces. In the Battle of San Juan Hill (actually Kettle Hill), Lt. Colonel Theodore Roosevelt earned a reputation as a military hero by leading the attack on entrenched Spanish positions. The regiment to which Roosevelt belonged, the First U.S. Volunteers, was recruited throughout the United States and known as the "Rough Riders" because of the large number of cowboys to volunteer. The 10th Cavalry, a regiment of black soldiers, supported the Rough Riders in the attack. Joseph Wheeler, a Confederate general of the Civil War, commanded U.S. forces in Cuba. Two of Robert E. Lee's nephews were also U.S. generals. The war ended eight months later with the signing of the Treaty of Paris on December 10, 1898. As a result Spain lost its control over the remains of its overseas empire. The treaty allowed the United States to purchase the Philippines Islands from Spain for $20 million. The war had cost the United States $250 million and 3,000 lives, of whom 90% had perished from infectious diseases. True to the letter of the Teller Amendment, American forces left Cuba in 1902. The Spanish-American War was seen domestically as a sign of increasing national unity. Hawaii. Kingdom of Hawaii. The Kingdom of Hawaii was established in 1795 with the subjugation of the smaller independent chiefdoms of Oʻahu, Maui, Molokaʻi, Lānaʻi, Kauaʻi and Niʻihau by the chiefdom of Hawaiʻi (or the "Big Island"), ruled by the dynasty of King Kamehameha the Great. In 1887 the Honolulu Rifle Company, a paramilitary force also known as the Honolulu Rifles, deposed the Hawaiian monarchy and forced King David Kalākaua to sign a new constitution at gunpoint. The bayonets fixed to their guns led to the term Bayonet Constitution. No voting rights were extended to Asiatics and the requirements for voting rights included land ownership. The Bayonet Constitution has become one of the most controversial documents in history. Native-born, European-descended Hawaiian Sanford B. Dole, serving as a friend of both Hawaiian royalty and the elite immigrant community, advocated the westernization of Hawaiian government and culture. Dole was a lawyer and jurist in the Hawaiian Islands as a kingdom, protectorate, republic and territory. King Kalākaua appointed Dole a justice of the Supreme Court of the Kingdom of Hawaii on December 28, 1887, and to a commission to revise judiciary laws on January 24, 1888. After Kalākaua's death, his sister Queen Liliʻuokalani appointed him to her Privy Council on August 31, 1891. Annexation. On January 17, 1893, the Queen, the last monarch of the Kingdom of Hawaiʻi, was deposed in a coup d'état led largely by American citizens opposed to her attempt to establish a new Constitution. Dole was named president of the Provisional Government of Hawaii formed after the coup, and was recognized within forty-eight hours by all nations with diplomatic ties to the Kingdom of Hawaii, with the exception of the United Kingdom. The Americans in Hawaii asked the US to annex the islands, but President Benjamin Harrison's annexation treaty was stalled in the Senate by Democrats until a Democratic President, Stephen Grover Cleveland, took office. With Grover Cleveland's election as President of the United States, the Provisional Government's hopes of annexation were derailed. In fact, Cleveland tried to directly help reinstate the monarchy, after an investigation led by James Henderson Blount. The Blount Report of July 17, 1893, commissioned by President Cleveland, concluded that the Committee of Safety conspired with U.S. ambassador John L. Stevens to land the United States Marine Corps, to forcibly remove Queen Liliʻuokalani from power, and declare a Provisional Government of Hawaii consisting of members from the Committee of Safety. Although unable to restore Lili'uokalani to her former position, Cleveland withdrew the treaty. The Territory of Hawaii or Hawaii Territory existed as a United States organized incorporated territory from July 7, 1898, until August 21, 1959, when its territory, with the exception of Johnston Atoll, was admitted to the Union as the fiftieth U.S. state. Resource Booms. There were a number of resource booms, resulting in the development of certain rural areas. Notable booms during this time include the Colorado Silver Boom, the Ohio Oil Rush, the Indiana gas boom, and the Cripple Creek Gold Rush. The Klondike Gold Rush was famous for showing there was value in Alaska with the discovery of Gold. 

Below are the authors of the XHTML tutorial: 

Progressivism. Industrialization led to the rise of big businesses at the expense of the worker. Factory laborers faced long hours, low wages, and unsanitary conditions. The large corporations protected themselves by allying with political parties. The parties, in turn, were controlled by party leaders, rather than by the voters. The Progressive movement was an effort to cure America: not so much an organized movement, but a general spirit of reform embraced by Americans with diverse goals and backgrounds during the early twentieth century . These problems included the aftermath of slavery, Reconstruction from the American Civil War, and female subjugation. The goal was to remove corrupt political machines and get more common people into the political process. Progressivism believed individuals could make the world better through regulation and reform. It worked on the Federal, state, local and public and private spheres, moving toward local public safety and efficiency, elimination of corruption, social justice, and the social control of knowledge. Local Reform. At the urban level, Progressivism mainly affected municipal government. The system whereby the city is governed by a powerful mayor and a council was replaced by the "council-manager" or the "commission" system. Under the council-manager system, the council would pass laws, while the manager would do no more than ensure their execution. The manager was essentially a weak mayor. Under the commission system, the executive would be composed of people who each controlled one area of government. The commission was essentially a multi-member, rather than single-member, executive. At the state level, several electoral reforms were made. Firstly, the secret ballot was introduced. Prior to the secret ballot, the ballots were colored papers printed by the political parties. Due to the lack of secrecy, bribing or blackmailing voters became common. It was to prevent businessmen or politicians from thus coercing voters that the secret ballot was introduced. Also, reforms were made to give voters more say in government. The initiative allowed voters to propose new laws. The referendum allowed certain laws (for example tax increases) to be approved by the voters first. Finally, the recall, allowed the voters to remove public officials for wrongdoing while in office. In addition, Progressives sought to combat the power of party leaders over which candidates would be nominated. The "direct primary" was instituted, under which the voters cast ballots to nominate candidates. Before the primary was introduced, the party leaders or party faithful were the only ones allowed to nominate candidates. The South pioneered some political reforms; "the direct primarily originated in North Carolina; the city commission plan arose in Galveston, Texas; and the city manager plan began in Stanton, Virgina. Progressive governors introduced business regulation, educational expansion,and other reforms that duplicated actions taken by northern counterparts. Labor Reforms. Progressive movement also attempted to give more power over legislation to the general populace. Three practices - the "referendum", the "initiative", and the "recall" - were created. The referendum allowed the voters to vote on a bill at an election before it took force as law. The initiative permitted the voters to petition and force the legislature to vote on a certain bill. Finally, the recall permitted voters to remove elected officials from office in the middle of the term. State laws were formed to improve labor conditions. Many states enacted factory inspection laws, and by 1916 nearly two-thirds of the states required compensation for the victims in industrial accidents. In 1901, Jane Addams founded the Juvenile Protective Association, a non-profit agency dedicated to protecting children from abuse. In 1903, Mary Harris Jones organized the Children's Crusade, a march of child workers from Kensington, Pennsylvania to the home of President Theodore Roosevelt in Oyster Bay, New York, bringing national attention to the issue of child labor. In 1909, President Roosevelt hosted the first White House Conference on Children, which continued to be held every decade through the 1970s. In 1912, the United States Children's Bureau was created in order to investigate "all matters pertaining to the welfare of children and child life among all classes of our people." At the instigation of middle class coalitions, many states enacted factory inspection laws, and by 1916 two-thirds of the states required compensation for victims of industrial accidents. An alliance of labor and humanitarian groups induced some legislatures to grant aid to mothers with dependent children. Under pressure from the National Child Labor to Committee, nearly every state set a minimum age for employment and limited hours that employers could make children work.Families that needed extra income evaded child labor restrictions by falsifying their children's ages to employers. States also regulated female labor by setting maximum work hours, especially when an accident at the Triangle Shirtwaist Factory resulted in the deaths of more than 100 women. The Supreme Court ruled in favor of regulated work hours for women in "Muller v. Oregon". Finally, some minimum wage provisions were introduced (for men and women). The Industrial Workers of the World (IWW) was founded in Chicago in 1905 at a convention of anarchist and socialist union members who were opposed to the policies of the American Federation of Labor (AFL). Unlike the AFL, which was a group composed of separate unions for each different trade (craft unionism), the IWW supported the concept of industrial unionism, in which all workers in a given industry are organized in one union, regardless of each worker's particular trade. They promoted the idea of "One Big Union" in the hopes that one large, centralized body would be better equipped to deal with similarly-large capitalist enterprises. 1881 The Great influx of Russian and Polish Jews. They formed an objectionable part of the population, because they couldn't speak English, lived closely crowded, and dirty. Penniless and unfamiliar with industrial conditions. They were apart of industrial, intellectual, and civil life. Their willingness to work 18 hours obnoxiously was crazy compared to Americans who worked (part-time). "It's not the condition that the immigrant comes from that determines he's usefulness;But the power one shows to rise above the condition." President Theodore Roosevelt. At the national level, Progressivism centered on defeating the power of large businesses. President Theodore Roosevelt, who succeeded to the Presidency when President McKinley was assassinated in 1901, helped the Progressive movement greatly. Coal Strike. In early 1902, anthracite coal miners struck. Their salaries had not been raised in over two decades. Furthermore, they were paid with "scrips". Scrips were essentially coupons for goods from pricey company stores. The president of the Reading Railroad, George F. Baer, said that the miners had erred by distrusting the owners. He declared that the mine owners were "Christian men of property to whom God has given control of the property rights of the country," who could be trusted more than union leaders. The owners and the miners refused to negotiate with each other. As autumn approached, many feared that the coal strike would cripple the economy. President Roosevelt intervened by asking the owners and miners to submit to arbitration. The miners accepted, but the owners refused Roosevelt's suggestion. Roosevelt then threatened to use the Army to take over the mines. The owners finally acquiesced; the strike was settled in 1903. Roosevelt's policy triumphed in 1904 when the Supreme Court, convinced by the government's arguments, created by J.P Morgan and his business allies. Roosevelt choose however , not to attack others trust, such as u.s steel another of Morgan's creations. Prosecution of northern securities began reportedly collared Roosevelt and offered "if we have done anything wrong, send your man to my man and we can fix it up". Sherman Antitrust Act. Roosevelt continued his Progressive actions when he revived the Sherman Antitrust Act. The Act sought to prevent companies from combining into trusts and gaining monopolies. A "trust" is formed when many companies loosely join together under a common board of directors to gain total control of an entire market so that prices can be raised without the threat of competitors. This total control of a market and subsequent price raising is a "monopoly". However, until Roosevelt's administration, the Act was rarely enforced. Hepburn Act. Roosevelt also enforced the Hepburn Act, which allowed the Interstate Commerce Commission to regulate railroads. The railroads had allied themselves with large businesses, charging higher rates to those business' competitors. Thus, the large businesses would gain even more power. The Hepburn Act prevented railroads from granting reduced rates to businesses. Panama Canal. President Roosevelt oversaw the successful completion of the Panama Canal. The finished canal vastly improved shipping logistics, allowing boats going from the Atlantic to the Pacific and vice versa to bypass the voyage around South America, saving much time on each trip. Conservation. Roosevelt also championed the cause of conservation. He set aside large amounts of land as part of the national park system. Conflicts with other Imperialist Nations. Imperialism was yet a common theme in the relations between nations in this era. It should be noted that although the US annexed Hawaii, Japan also had interests in the island and an aggressive foreign policy; Japan had already seized Taiwan from China in 1885 and would annex Korea in 1905. Imperial Germany was another aggressive power. The U.S. and Germany had conflicts over who would control Samoa, in the Pacific, as well as nearly faced a naval war with Germany in 1902 over German plans to seize the customs revenues of Venezuela. However, under the administration of President Theodore Roosevelt, the United States became more open in asserting international power. In 1905, Russia and Japan went to war over control of Korea and China. The Japanese won naval victories over two Russian fleets, in the Battles of the Yellow Sea and Tsushima. President Roosevelt offered to negotiate peace between the two nations, and in Portsmouth, New Hampshire, a peace treaty was signed. To demonstrate the ability of the United States to project power around the world (unlike Russia), Roosevelt ordered a fleet of U.S. men-o'war to sail around the world. The fleet left the east coast of the U.S. in 1908 and returned in 1909, visiting ports in Europe, Australia, and Japan. President William Howard Taft. When Theodore Roosevelt decided not to run for the presidency again in the election of 1908, that opened the doors for another Republican candidate. Into the gap stepped William H. Taft, former Ohio Supreme Court Justice, former Solicitor General of the United States, and former judge on the United States Court of Appeals for the Sixth Circuit. James S. Sherman was chosen as his vice president. They ran against the Democratic slate of William Jennings Bryan and John Kern, and Socialist nominee Eugene V. Debs. Taft won the election by over a million votes and the Republicans retained control of both houses of Congress. During Taft's presidency, his primary goals were to continue Roosevelt's trust-busting and to reconcile old guard conservatives and young progressive reformers in the Republican Party. Taft was somewhat more cautious and quiet than Roosevelt, despite his own credentials, and therefore had less public attention. Taft tried to win over the Filipino people by reforming education, transportation, and health care. New railroads, bridges, and telegraph lines strengthened the economy. A public school system was founded, and new health care policies virtually eliminated such diseases as cholera and smallpox. These reforms slowly reduced Filipino hostility. Although Taft was less of an attention-grabber than Roosevelt, he went far beyond what Roosevelt ever did. Taft used the Sherman Antitrust Act, a law passed in 1890 that made trusts and monopolies illegal, and they had to sue many large and economically damaging corporations. Taft won more antitrust lawsuits in four years than Roosevelt had won in seven. Taft also pushed for the Sixteenth Amendment, which gave the federal government the right to tax citizens' income. The amendment was meant to supply the government with cash to replace the revenue generated from tariffs, which Progressives had hoped that Taft would lower. Taft failed to lower the tariff. In addition, he failed to fight for conservation and environmentalism, actually weakening some conservation policies to favor business. When Roosevelt came back from an expedition to Africa in 1910, he was disappointed in Taft, and vigorously campaigned for progressive republicans in the congressional elections of 1910. Roosevelt tried to capitalize on his still enormous popularity by again running for reelection in 1912, but he failed to win the nomination because of Taft's connections to influential people in the Republican Party. Roosevelt and his supporters then broke off from the Republicans, forming the Progressive Party. (This later was known as the Bull Moose party after Roosevelt declared that he felt "as strong as a bull moose!") The Republican split hurt the two candidates, and Democratic candidate Woodrow Wilson gathered a 42 percent plurality of the popular votes and 435 out of 531 electoral votes. On the Supreme Court. Taft would later become a member of the Supreme Court, making him the only former President to do so. In 1921, when Chief Justice Edward Douglass White died, President Warren G. Harding nominated Taft to take his place, thereby fulfilling Taft's lifelong ambition to become Chief Justice of the United States. Very little opposition existed to the nomination, and the Senate approved him 60-4 in a secret session, but the roll call of the vote has never been made public. He readily took up the position, serving until 1930. As such, he became the only President to serve as Chief Justice, and thus is also the only former President to swear in subsequent Presidents, giving the oath of office to both Calvin Coolidge (in 1925) and Herbert Hoover (in 1929). He remains the only person to have led both the Executive and Judicial branches of the United States government. He considered his time as Chief Justice to be the highest point of his career: he allegedly once remarked, "I don't remember that I ever was President." President Woodrow Wilson. Although Woodrow Wilson was a Democrat, he still pushed for progressive reforms. One of the first successes of his administration was the lowering of tariffs, which he accomplished in 1913. Wilson believed that increased foreign competition would spur U.S. based manufacturers to lower prices and improve their goods. That same year, Wilson passed the Federal Reserve Act, which created twelve regional banks that would be run by a central board in the capitol. This system gave the government more control over banking activities. Wilson also pushed for governmental control over business. In 1914, a Democratic-controlled Congress established the Federal Trade Commission (FTC) to investigate companies that participated in suspected unfair and illegal trade practices. Wilson also supported the Clayton Antitrust Act, which joined the Sherman Antitrust Act as one the government's tools to fight trusts the same year. By the end of Wilson's First term, progressives had won many victories. The entire movement lost steam, though, as Americans became much more interested in international affairs, especially the war that had broken out in Europe in 1914. The Supreme Court and Labor. Upset workers had succeeded in lobbying Congress to pass legislation that improved work conditions. However, the Supreme Court of the United States somewhat limited the range of these acts. In Holden v. Hardy (1896), the Supreme Court ruled that miners' hours must be short because long hours made the job too dangerous. However, in Lochner v. New York, laws ruled that bakery workers did not have a job dangerous enough to put restrictions on the free sale of labor. Putting aside this decision, in 1908, the decision in Muller v. Oregon said that women's" health must be protected "to preserve the strength and vigor of the race." This did," clearly, protect women's health, but it also locked them into menial jobs. Controlling Prostitution. Moral outrage erupted when muckraking journalists charged that international gangs were kidnapping young women and forcing them into prostitution, a practice called white slavery. Accusations were exaggerated, but they alarmed some moralists who falsely perceived a link between immigration and prostitution. Although some women voluntarily entered "the profession" because it offered income and independence from their male counterparts, some women had very little option to a life where they had little if any amenities and many were forced into this profession and lifestyle. Reformers nonetheless believed they could attack prostitution by punishing both those who promoted it and those who practiced it. In 1910 Congress passed the White Slave Traffic Act (Mann Act), prohibiting interstate and international transportation of a woman for immoral purposes. By 1915 nearly all states outlawed brothels and solicitation of sex. Such laws ostensibly protected young women from exploitation, but in reality they failed to address the more serious problem of sexual violence that women suffered at the hands of family members, presumed friends, and employers. Football and the Formation of the NCAA. By the turn of the century American football was already in the process of becoming a large national sport. Originally formed and played at universities as an intercollegiate sport, it was seen as only for the upper class. The size of the field depended on what the players agreed with, but it was almost always over 100 yards. Once a player started a game, the player could not leave unless he/she became injured. Very soon the sport began to gain spectators, and with spectators came controversy. With over 15 deaths in 1905 alone, many saw a need for change in the sport. However, others liked the violence and would watch because of this. President Roosevelt formed a group to reconstruct the rules of football and make it less violent. Standard rules would not be made and used until 1894. The group was originally named the Intercollegiate Athletic Association, and in 1910 it was renamed to the National College Athletic Association. 

Politics and Government. Presidency of Warren G. Harding. A new sympathy toward business was shown in the election of Republican Warren G. Harding as president in 1920. His administration helped streamline federal spending with the Budgeting and Accounting Act of 1921, supported anti-lynching legislation (which was, however, rejected by Congress), and approved bills assisting farm cooperatives and liberalizing farm credit. Scandals. The Harding administration was also known for its scandals. He had had an affair with the wife of an Ohio merchant: the resulting daughter was never officially acknowledged. He also appointed some cronies, who saw office as an invitation to personal gain. One of those men was Charles Forbes, head of the Veterans Bureau. He was convicted of fraud and bribery in connection with government contracts, and was sent to prison. Another crony, Attorney General Harry Daugherty, was involved with an illegal liquor scheme. He only escaped prosecution by refusing to testify against himself. Teapot Dome Scandal. The most notorious of these scandals was the revelation that the Secretary of the Interior, Albert Fall, accepted bribes to lease government property to private oil companies in the Teapot Dome Scandal. The popular conservation legislation created by Harding's predecessors, presidents Teddy Roosevelt, William Taft, and Woodrow Wilson, had set aside naval petroleum reserves in Wyoming and California. Three naval oil fields, Elk Hills and Buena Vista Hills in California and Teapot Dome in Wyoming, were tracts of public land meant as emergency underground supplies to be used by the navy only when regular oil supplies diminished. Teapot Dome received its name because of a rock resembling a teapot above the oil-bearing land. Politicians and private oil interests had opposed the restrictions placed on the oil fields, claiming that the reserves were unnecessary, and that the American oil companies could provide for the U.S. Navy. Civil and criminal suits concerning Teapot Dome lasted through the 1920s. In 1927 the Supreme Court finally ruled that the oil leases had been corruptly obtained and invalidated the Elk Hills lease in February of that year and the Teapot lease in October of the same year. The navy regained control of Teapot Dome and Elk Hills reserves. Albert Fall was found guilty of bribery in 1929, fined $100,000, and sentenced to one year in prison. Harry Sinclair refused to cooperate with the government investigators, was charged with contempt, and received a short sentence for jury tampering. Edward Doheny was acquitted in 1930 of attempts to bribe Fall. The Teapot Dome scandal was a victory for neither political party. It became a major issue in the presidential election of 1924, but neither party could claim full credit for divulging the wrongdoing. It became the first true evidence of government corruption in America. The scandal revealed the problem of natural resource scarcity and the need to protect for the future against emergency depletion of resources. Vice President Calvin Coolidge, who assumed the presidency after Harding's death, handled the problem very systematically, and his administration avoided any damage to its reputation. Technology. Although there were innumerable technical innovations, the vast changes in American life about this time had two major technical bases, mass production (the assembly line), and mass testing. In the vast steel factories and in cloth mills individuals had to move together with the machines. Any mistake could lead to an accident, perhaps a fatal one. Henry Ford's assembly line worked on the same principle, but went much further. A car went from station to station, from worker to worker. Each worker had one function—tightening nuts, adding a component—and only that function, as if he were himself a machine. His main interest was in doing those motions which would do his job and do it most efficiently. (In this respect the system drew upon work efficiency experts such as Frederick Winslow Taylor and the Gilbreths.) The advantages to this extreme systemization were fundamental. As with the cloth factories, the product was produced extremely quickly at all hours of the day or night. Very little training was needed for those jobs. The results of this system were extremely long-lasting. The Model T was seen as a durable car, and "the tin Lizzie" retained public affection even when it was superseded by cars with self-starters. The cars were also affordable, with results as seen below. Henry Ford raised his wages regularly, urging that the men who made the cars also buy them. A small demerit was that these new cars were extremely ugly. The Stanley Steamer had been sleek, with lines like what would later be called "streamlining." Dusenbergs and Pierce-Arrows had a variety of hues and such extras as bud vases. Ford famously said that his purchasers could have any color they wanted, "so long as it is black." As people became more prosperous, they could shop for colored paint jobs and detailing for their luxury cars. However, creating affordable and beautiful goods was a movement away from Ford's version of the market. More importantly, working on the assembly line wore on the workers. Standing in one place and squinting, working with a few muscles for hours a day (or night) could be very fatiguing. Human beings weren't made to live like that. When there were a limited set of priorities for working, workers could be easily replaced, just like machines. Not every boss of the assembly line paid as well as Ford. Mass testing was a requirement for the assembly line—a bad part made a bad product. But it had actually begun as a human policy, in the requirements for the late 19th Century census. It was accelerated in the desire to find sound men for the First World War. Psychologists were employed to create intelligence tests to weed out unfit soldiers. Their weapons had to be carefully inspected, for when shoddy goods reached the front lines, the result could be a disaster. In the post-war world, Big Business began developing research and development departments. Before a change was implemented, there had to be a prototype, and the effects on the public had to be carefully measured. Economics became a matter, not merely of becoming prosperous, but of selling to the largest number possible. (The term "mass market" originated in the 1920s.) The Automobile. In the 1920s, the United States automobile industry began an extraordinary period of growth by means of the assembly line in manufacturing. Cars began to alter the American lifestyle. In 1929, one out of every five Americans had a car. They began using their own automobiles instead of the street cars. Cars also replaced horses. This made the streets cleaner, because there wasn't as much horse manure. (However, this was replaced by other, more subtle forms of pollution. In the 1920s gasoline companies started adding lead to their fuel to increase engine efficiency.) The idea of "homes on wheels" was also created around this time. Americans were packing up food and camping equipment in order to get away from home. By the 1920s most automobiles gained cloth or steel roofs, offering a private space for courtship and sex. Women gained from the automobile revolution. Women who learned to drive achieved new-found independence, taking touring trips with female friends, conquering muddy roads, and making repairs when their vehicles broke down. Prosperous African Americans for the first time obtained a limited freedom from local discrimination. A family could drive around and past "closed" White communities, and to beaches, camps, and other holiday destinations. (However, the family car would have to carry its own food, drink and gas, and not stop before it reached its destination. The largest-scale pamphlet for "safe" businesses African Americans could use, the Negro Motorist Green Book, was only published beginning in 1936.) The car was the ultimate social equalizer. There were 108 automobile manufacturers in 1923 and colors allowed owners to express personal tastes. An abundance of fuel fed these cars. In 1920, the United States produced sixty-five percent of the World's oil. Road construction was extensive. The first timed stop-and-go traffic light was in 1924. Industries related to the manufacturing and use of automobiles also grew; petroleum, steel, and glass were in high demand, leading to growth and profitability in related sectors. State governments began to build roads and highways in rural areas. Gasoline stations were installed across the country, evidence of the sudden and continued growth of the petroleum industry. Automobile dealers introduced the installment plan, a financing concept that was adopted in many other parts of business. Thus, the automobile industry's growth had repercussions throughout the nation. With a perfected design of Henry Ford's assembly line automobiles began to be more affordable for the common US citizens all over the country. A lot of men were hired to work in car factories. Health and Life Expectancy. The relation between food and health had long been known. For example, since the 18th century it has been known how to fight scurvy, and mariners have taken fruit on long voyages. Yet the fact that scurvy is caused by lack of vitamin c was only discovered in 1932. From 1915 to the end of the 1920s most vitamins were discovered. Food regulation began to ensure a safer food supply. People began to have access to and the possibility of choosing more and better food, due to faster transport and refrigeration. Technical information was also more easily transmitted, and by 1930 nutritionists began to emphasize to the public the need for consumption of certain foods, and their constituent vitamins and minerals, on a daily basis. Food companies began marketing their products, on how their products contain certain amounts of your daily vitamins and therefore healthy. However, the advertisements sometimes contained unusual ideas about nutrition. For example, some candy bars were advertised by their "food value." And Welch's Grape Juice marketed their product as containing nutrients and vitamins, but failed to inform the reader of the large amount of sugar also included. But the emphasis on nutrition and good hygiene made many Americans healthier. This was the decade when penicillin and insulin were discovered. During this time the life expectancy at birth in the United States also increased from fifty-four to sixty percent, and infant mortality rate decreased by one-third. However this was not the case for nonwhites: the mortality rate for nonWhite children was about fifty to one hundred times that of Whites during this era. (Rickets among the poor and among rural African Americans was seen as the result of poor genetics, "bad blood." The American fad for Eugenics and the sterilization movement also grew in this era.) Accident fatalities also increased by roughly 150 percent, for the car was becoming faster and more common. Elderly Americans and Retirement. As more adults survived into old age, an interest in pensions and other forms of old age assistance grew. In the third decade of the 20th century, one third of Americans sixty-five and older depended financially on someone else. Over the past fifty years many European countries had established state-supported pension systems. In 1923 the Pennsylvania Chamber of Commerce called old-age assistance “un-American &amp; socialistic.” But during the 1920s state resistance to pension plans eroded. Isaac Max Rubinow and Abraham Epstein attempted to persuade legislators and associations such as labor unions to endorse old-age assistance. By 1933 almost every state gave some minimal support to needy elderly. Culture of the 1920's. During this time period, new social values emerged. It became difficult to determine what was socially acceptable, as youth frequently took up smoking, drinking, and a new openness about sex. They were being influenced less by their parents and more by their peers and schoolmates. Schools in the cities geared up for mass education, segregating children with others of their same age. The rite of passage, dating without adult supervision, became more commonplace among these youth. The Flapper was the female symbol of this change, as the raccoon-coated Sheik was the symbol among young men. The dresses then in fashion de-emphasized the bodice, with a flat abdomen, the so-called "boyish figure." Flappers did not have the long hair of their mothers and grandmothers, but short, "bobbed" styles. They drank and smoked like men, knew all the latest dances and songs, and openly swore and talked of sex. It is unknown how frequent the Flapper really was. Bobbed hair was fashionable among women and girls, but there was never a standard measurement of this or any related trend. Movie actresses such as Louise Brooks and Clara Bow were shown drinking on the big screen, and F. Scott Fitzgerald's writings showed a literary interest in the Flapper. During the War, servicemen became used to lectures on preventing venereal disease, and thus became more comfortable with the idea of contraception. Condoms started to be made with latex instead of animal tissues, and became a product which could be mass-produced on an assembly line. Birth control became more available, and more respectable. With a greater chance for babies to survive infancy, and with the ability to time when they came into the world, the number of children in a middle-class family began to go from four or five to two or three. Unfortunately, this overlapped with the eugenics movement. As the middle class became more mobile, it was much less able to rely on the advice of grandparents and family, and "expert" child care advice became popular. This advice was different from that commonly used nowadays. John B. Watson, who published his book in 1928, advised against picking up infants, holding them when they cried, or cossetting them or showing them too much affection. Radio. Radio had been used for ship-to-shore communication since the Titanic sent out a Morse code S-O-S. It was used by both sides during the First World War. Wilson considered nationalizing the medium, as Great Britain was later to nationalize the British Broadcasting Corporation, but corporate outcry overruled him. By 1920, thousands of curious machines produced screeches for the hobbyist, with an occasional, distant snatch of voices or music. In 1920 the assembly line did its work, producing an RCA "Cat's Whisker" receiver for under four dollars. In October of that year Westinghouse created the first radio station, KDKA. In November it provided running coverage of the Presidential elections. By the mid-'20s programming ran from morning till night. In 1924, the first radio network, the National Broadcasting Company, began operations between New York and Boston. In 1927, the Columbia Broadcasting System began. The Federal Radio Commission was set up in 1926, and organized in the Radio Act of 1927. Advertisers sponsored programs: one popular music program was The A &amp; P Gypsies, giving coast-to-coast coverage to the A &amp; P grocery stores. The news and entertainment provided was vetted by the sponsor, and anything which would offend sponsors was forbidden. But within those lines much was permitted. One could occasionally find high culture (though the Metropolitan Opera broadcasts only began at the end of 1931), but the aim was to air songs in the middle range of culture; "The Lost Chord," or "Drink to Me Only With Thine Eyes." They played popular music, but not much jazz: Paul Whiteman, the so-called "King of Jazz," was not. (There were exceptions to this rule; some high-powered radio stations in Mexico poured out jazz, "Black music," and ads for toxic patent medicines.) Much of the country's culture was not covered, though the Grand Ole Opry began its broadcasts in 1925. But as electrification expanded, the market for radio grew, and some stations experimented. A pair of White comedians, Freeman Gosden and Charles Correll, created a comic, sentimental serial drama, "Amos 'n' Andy". At a time when lynchings of African Americans occurred as far North as Ohio, this was a comedy about two stupid African Americans who mispronounced their words. But it also created sympathy for them. "One episode ended with Amos and Andy in desperate need of a typewriter; nearly two thousand typewriters were immediately sent in by listeners." Yet " "Amos 'n' Andy" 's popularity was no doubt due to excitement over this new national experience. For the first time Americans could all enjoy the same event at the same moment." Movies. In the 1920s movies also grew into a popular recreation. By 1922, about 40 million people were going to the movies each week; that number jumped to about 100 million people by the end of the decade. Movie stars such as Douglas Fairbanks, Mary Pickford, and Charlie Chaplin became known around the world. Eight studios dominated the industry, consolidating and integrating all aspects of a film's development. By 1929, the film-making firms that were to rule and monopolize Hollywood for the next half-century were the giants or the majors, sometimes dubbed The Big Five. The Big Five studios were Warner Bros., RKO, Paramount, Metro-Goldwyn-Mayer, and Fox Film Corporation. They produced more than 90 percent of the fiction films in America and distributed their films both nationally and internationally. Each studio somewhat differentiated its products from other studios. A movie house was only allowed to play the products of one studio. Thus, for example, the New York Paramount only played cartoons, newsreels, and fiction films created by Paramount Studios. (In the 1920s Paramount distributed the work of Max Fleischer Studios, creator of the Koko the Clown Cartoons.) Each division of the studio was contracted to make so many films each year. If a movie house wanted to get the films of a Gloria Swanson or a Rudolf Valentino, it had to accept a given number of films by a less-liked star. This "block booking" ensured that certain actors got publicity and kept the screens under the thumb of the studio. However, in return each theater was ensured of a weekly change of movies, with the full backing of the studio. In addition to the projectionist, ushers and candy and cigarette sellers, the Paramount Theater employed a grand musician to accompany the silent film on one of the largest theater organs ever created. Its halls were ornamented by hand-painted murals. The top-line theaters were called "movie palaces." The most popular studio movies often used sweepingly romantic stories set in exotic lands: Argentina in 1921's "The Four Horsemen of the Apocalypse", Persia in 1924's "The Thief of Bagdad". The 1925 movie "Ben Hur" was shot partially in Italy and partially on huge purpose-built sets in California. It had 42 cameras shooting the still-famous chariot race. Among the famous or yet-to-be-famous figures swelling the scenes as extras were the Barrymore brothers, Lillian and Dorothy Gish, comedian Harold Lloyd, William Randolph Hearst's love Marion Davies, and studio head Samuel Goldwyn. The religious sequences used two-tone technicolor. It was the most expensive film yet made. Prohibition. Although total alcohol consumption halved, some people blatantly disregarded Prohibition. There were loopholes in the Volstead Act, the twenty-two page law which defined Prohibition. Churches could use wine in their ceremonies, and alcohol drunk as a medicine (this was still part of the medical profession) was still allowed. The amount of "religious" and "medicinal" wine suddenly increased. Some illegal alcohol was imported from Canada, Cuba, and Mexico, which never made alcohol illegal. Some was home-made American. "Bootleggers" were found in many places throughout the country, from backwoods stills (illegal alcohol production had continued after the Whiskey Rebellion) to urban "bathtub gin." The Volstead Act had said that personal consumption of alcohol in one's own home was legal, though it had prohibited public gatherings to drink. The occasional secret saloons called "speakeasies" which sprang up in cities were therefore illegal. These required money, and a new criminal underworld rose to fund them and profit from them. Some of this money funded pay-offs to police to stop enforcement of Prohibition. Gangs prospered in this hidden economy. Many jobs came out of Prohibition, both from alcohol and from the "front" legitimate businesses set up to launder speakeasy money. However, these jobs came with great risks, from blackmail and graft to outright violence. Some commentators felt that Prohibition was too harsh and that it made a criminal out of the average American man or woman, who would have bought alcohol legally if it were available. Gangs and Violence. There was obviously a huge market for what in the 1920s was an illegal commodity. Gangsters provided this commodity. Major gangsters in this period included Charles "Lucky" Luciano, Mayer Lansky, and "Dutch" Schultz. Perhaps the most notorious was Chicago's Al Capone. Capone smuggled alcohol all over the Midwest. He was also responsible for drug smuggling and murder, and bribed both police and important politicians. Despite the deference given Capone by "bought" figures, he had enemies from other Chicago gangs. He rode in an armor-plated limousine, always accompanied by armed bodyguards. Violence was a daily occurrence in Chicago. 227 gangsters were killed in the space of four years. On St Valentine's Day, 1929, seven members of the O'Banion gang were shot dead by gangsters dressed as police officers. In 1931, the government got around the corrupted regular police by arresting Capone for tax evasion, rather than for his many violent offenses. He got eleven years in jail, and left prison with his health broken. Bonnie and Clyde. Bonnie and Clyde were also a famous pair of murderers and thieves in the 1920s during the prohibition era with their gang. Clyde Champion Barrow and his companion, Bonnie Parker, were shot to death by officers in an ambush near Sailes, Bienville Parish, Louisiana on May 23, 1934, after one of the most colorful and spectacular manhunts the nation had seen up to that time. Barrow was suspected of numerous killings and was wanted for murder, robbery, and state charges of kidnapping. The Federal Bureau of Investigation (FBI), then called the Bureau of Investigation, became interested in Barrow and his paramour late in December 1932 through a singular bit of evidence. A Ford automobile, which had been stolen in Pawhuska, Oklahoma, was found abandoned near Jackson, Michigan in September of that year. At Pawhuska, it was learned another Ford car had been abandoned there which had been stolen in Illinois. A search of this car revealed it had been occupied by a man and a woman, indicated by abandoned articles therein. In this car was found a prescription bottle, which led special agents to a drug store in Nacogdoches, Texas, where investigation disclosed the woman for whom the prescription had been filled was Clyde Barrow's aunt. Further investigation revealed that the woman who obtained the prescription had been visited recently by Clyde Barrow, Bonnie Parker, and Clyde's brother, L. C. Barrow. It also was learned that these three were driving a Ford car, identified as the one stolen in Illinois. It was further shown that L. C. Barrow had secured the empty prescription bottle from a son of the woman who had originally obtained it. On May 20, 1933, the United States Commissioner at Dallas, Texas, issued a warrant against Clyde Barrow and Bonnie Parker, charging them with the interstate transportation, from Dallas, to Oklahoma, of the automobile stolen in Illinois. The FBI then started its hunt for this elusive pair. Religions and Revivalism. Just as the Religious section of the newspaper had long been popular, the new medium of radio became a way to increase religious visibility. Churches bought broadcasting slots from stations eager to seem like good neighbors. City stations might broadcast programs of interest to Catholics and Jews, as well as from minority faiths or cults. This impinged on listeners, coming into their homes, as printed media did not. The religious revival had been a feature of both mainline denominations and smaller sects since the turn of the century. Both mainline and independent preachers called upon listeners to give up frivolity and turn to a purer faith. Many of these sermons condemned movies and theatre, novels and card gambling, drinking and modern fashion, including women's short dresses and makeup. Many of them had supported the imposition of Prohibition. The Twenties provided electric lighting, amplification, and radio coverage for revivals. Some popular preachers traveled by train or motor car to cities and towns across the country. Aimee Semple McPherson and Billy Sunday were among the most notable of these, and aroused controversy. Jazz. Jazz is an American musical art form which originated around the beginning of the 20th century in Black communities in the Southern United States from a confluence of African and European music traditions. The “hometown” of jazz is considered to be in New Orleans. Early jazz musicians would called New Orleans their home even if they have never been there. Jazz employed a number of Black men and women. Jazz spread through America very quickly. The style's West African pedigree is evident in its use of blue notes, call-and-response, improvisation, polyrhythms, syncopation, and the swung note of ragtime. Beginning in 1922, Gennett Records began recording jazz groups performing in Chicago. The first group they recorded was the New Orleans Rhythm Kings, followed in 1923 by King Oliver's Creole Jazz Band with Louis Armstrong. Another indie company in Chicago, Paramount Records, was competing with Gennett and Okeh for jazz talent. The Black community took notice: authors such as Langston Hughes often mentioned the music in his poems, both positively and negatively. Business Overseas. After the war, many manufacturing companies faced hard times as they attempted to convert from wartime production of weapons and planes to what they had traditionally produced before the war. However, the pro-business policies put in place first by Harding, then Coolidge, allowed business to flourish. While business did well at home—the raising of tariff rates from 27% (under the Underwood-Simmons Tariff) to 41% certainly helped in this regard—many major companies did quite well overseas. Just as these companies had started to do before the war, they set up shop in a variety of countries based around the resources located there. Meat packers like Gustavus Swift went to Argentina; fruit growers went to Costa Rica, Honduras, and Guatemala; sugar plantation owners went to Cuba; rubber plantation owners to the Philippines, Sumatra, and Malaya; copper corporations to Chile; and oil companies to Mexico and Venezuela (which remains today a great source for oil). Steamships and telegraphs made for easy transport and communication. Organized Labor. The Organized labor force during the 1920s suffered a great deal. During this time the country was fearful of the spread of communism in America, because of this widespread fear public opinion was against any worker who attempted to disrupt the order of the working class. The public was so anti-labor union that in 1922 the Harding administration was able to get a court injunction to destroy a railroad workers strike that was about 400,000 strong. Also in 1922 the government took part in putting to an end a nationwide miners strike that consisted of about 650,000 miners. The federal and state level of government had no toleration for strikes, and allowed for businesses to sue the unions for any damages done during a strike. Major Cases. The Sacco-Vanzetti Trial, Leopold and Loeb, the Scopes Trial, and the Black Sox Trial were all significant court cases during the 1920s. Each of these court cases were unique and monumental in their own right, and set a precedent for the years to come. The Scopes Trial. In 1925, John Thomas Scopes, a biology teacher, was tried and convicted in Tennessee for teaching about evolution in his public school classroom as an explanation of the origin of humans, as opposed to the Biblical story of Adam and Eve, which was supported by state law at the time. This was a major dispute and caught the attention of many popular government officials such as William Jennings Bryan, who spoke on behalf of the prosecution. Bryan saw evolution as not only pernicious in its own right, but as a platform for eugenics: the textbook that Scopes used, "Civic Biology", advocated "racial hygiene." Although the modernists lost the case, they still were happy to have highlighted the illogical reasoning behind the law that schools could not teach alternative theories for the origin of man. They were also happy that this trial and conviction didn't affect the expansion of fundamentalist ideals. The Southern Baptist Convention, a Protestant group, became one of the fastest growing denominations after the trial showing that it may have even given popularity to the religious denomination. The beliefs of these groups resulted in an independent subculture with their own schools, radio programs, and missionary societies. Minority Women. During the 1920s there were almost double the amount of minority women than White in the workforce. Women, especially minorities, who held factory jobs held the least desirable and lowest paying jobs in factories. Black women mostly held domestic jobs such as cooking, cleaning, and laundry. There were many openings for educated minority women in the social work, teaching, and nursing fields during this time, however they faced much discrimination. The economic needs of the family brought thousands of minority women into having to work. Mexican women, mainly in the Southwest worked as domestic servants, operatives in garment factories, and as agricultural laborers. This was looked down upon because the Mexican culture traditionally was against women labor. Next to African women, Japanese women were the most likely to hold low paying jobs in the work force, they worked in the lowest paying jobs; they faced very strong racial biases and discrimination on a regular basis as well. African-Americans and the Ku Klux Klan. Rebirth of the Ku Klux Klan. Southern states segregated public facilities (like buses). In half the South fewer than 10% of the Blacks were allowed to vote. The Ku Klux Klan flourished 1921–26 with a membership of millions of Protestants. Not only was the Ku Klux Klan big in the south, but in such northern states as Ohio, Oregon, and Indiana. Indiana's governor and an Oregon mayor were both members of the KKK. Many KKK members were women, nearly a half million in women's auxiliary associations. Klansmen organized marches and violence against African-Americans, Catholics, and Jews, as well as bootleggers and adulterers. They gained new support from nativists who had detested the mass immigration to the Northeast in the early 1900s. The return of the Klan caused a split in the Democratic Party which allowed Calvin Coolidge, a conservative Republican, to take office in 1924. Blacks were widely persecuted by the Ku Klux Klan, but they were not the only group of people that the KKK targeted because they believed in “Native, White, Protestant supremacy.” They also targeted groups like Mexicans, Jews and Catholics. A Klu Klux Klan newspaper ran a doggerel poem of a dialogue between the Pope and the Devil, with the latter saying the KKK "will make it hotter than I can for you in hell." The Ku Klux Klan would also try and bring justice into their own hands when it came to dealing with bootleggers, wife beaters and adulterers, and even the Knights of Columbus. The Great Migration. In most cities, the only way Blacks could relieve the pressure of crowding that resulted from increasing migration was to expand residential borders into surrounding previously White neighborhoods, a process that often resulted in harassment and attacked by White residents whose intolerant attitudes were intensified by fears that Black neighbors would cause property values to decline. Moreover, the increased presence of Blacks in cities, North and South, as well as their competition with Whites for housing, jobs, and political influence sparked a series of race riots. The Tulsa Race Riots of 1921 and the Rosewood Massacre of 1923 (where a town was wiped off the map by racist violence) were examples of this extreme of hatred. The Tulsa Race Riots involved the local white community looting a prosperous and thriving black community, and destroying it with arson and an aerial bombing by plane. The Great Migration ignited hatred toward Blacks in Northern big cities. Blacks migrated to cities such as New York, Chicago, and Detroit, and in Western cities such as Los Angeles and San Diego. Although most of the migrants were poor and lived in cheap urban housing, some were able to afford better houses in White neighborhoods. However, even prosperous people were unable to live where they wanted. Discrimination could be as open as the notice in Wanted ads -- "No Negroes allowed"—or as quiet as the refusal of a real estate agent. The brave Black family who actually bought in those neighborhoods would face snubs from the neighbors, refusal of services from local businesses, and sometimes covert or open violence. (An example of this last is seen in Detroit's Ossian Sweet case of 1925.) Organizing. To fight this discrimination many Black movement groups formed. The Universal Negro Improvement Association (UNIA) was formed by Marcus Garvey, an immigrant from Jamaica living in Harlem. Garvey preached a message of equality that many, including other Black leaders, considered radical. Garvey helped start companies and news papers directed towards the Black community. He gained many followers around the US, especially in cities. Amritjit Singh estimates that Garvey and the UNIA had over half a million followers. Garvey created more racial cohesion and inspired the Black community to stand up. Yet others, including the prominent Black author W.E.B. Du Bois, considered Garvey's approach extreme and believed that it would backfire. Du Bois believed in a "gradualist" strategy, working through education and the legal system. He and some other Black leaders petitioned the U.S. Attorney General and had Garvey deported back to Jamaica. Yet Garvey's message lived on long after he was deported, and he was one of the early inspirations of 1960 civil rights leader Malcolm X. The Harlem Renaissance. In New York City's Harlem and in half a dozen other Northern cities, a Black culture began to form as a result of the relative economic and educational advantages given through the Great Migration. Black businesses, legal and political systems, and arts societies flourished. Here "the talented tenth" had its own business and fraternal institutions. Poets Langston Hughes, Claude McKay, and Jean Toomer, and musicians such as Chick Webb and Duke Ellington were published in mainstream magazines and heard in White-frequented (though sometimes segregated) clubs. The Harlem Renaissance talked of the contemporary Black community's hopes and fears. Other races during the 1920's. The Indian Citizenship Act of 1924 gave indigenous people in the United States citizenship. The Immigration Act of 1924 effectively eliminated immigration from Asia, and limited the immigration of Jews, Italians, and Eastern Europeans. The End of Prosperity and the Stock Market Crash of 1929. In the 1920s, farmers did not do so well. A lot of farms did not have running water or electricity, and pay was low due to surplus. World War I had disrupted farming in Europe and the warring European nations greatly depended on American farming for food. When peace came, demand for crops like cotton and grain suddenly fell but farmers kept planting at wartime rates, so they were left without money to pay off their loans or new devices like tractors. A lot of farmers were dependent growing cotton. However, in the twenties the price of cotton plummeted because of new man-made materials that entered the market. Matters were made worse by the invasion of the boll weevil, an insect which planted its eggs in the boll (cotton blossom), and ate the cotton. The Southern economy was partially saved through following the urging of inventor George Washington Carver and planting peanuts instead of cotton. In 1925-1927 George Washington Carver patented two uses for peanuts, and hundreds of more inventions from soybeans, pecans, and even sweet potatoes. Some inventions he made from peanuts and soybeans are paper, instant coffee, shaving cream, mayonnaise, soap, and talcum powder. None of these procedures were ever recorded by him in a notebook. He urged increased participation of Blacks in agricultural education. On October 24, 1929, today known as Black Thursday, the stock market began its downhill drop. After the first hour, the prices had gone down at an amazing speed. Some people thought that after that day, the prices would rise again just as it had done before. But prices kept dropping. On October 29, 1929, Black Tuesday, more than 16 million shares were sold, but by the end of the day, most stocks ended below their previous value, and some stocks became totally worthless. By November 13, the prices had hit rock bottom. The stock AT&amp;T had gone from 304 dollars to 197. Much of America had celebrated unheard of prosperity for eight years, but the Stock Market Crash put an end to that within a few weeks. Questions For Review. 1. Name the economic effects of one of the following: the automobile; mass production as a whole; the boll weevil. ← World War I · US History · Great Depression and New Deal → 

Please see COSTP World History Project for the parallel project based on California content standards (note that it has no content). Newly Adopted AP World History Standards. World History/AP World History Standard Previous Standards. California standards. World History/California Content Standard http://www.cde.ca.gov/re/pn/fd/documents/hist-social-sci-frame.pdf http://www.cde.ca.gov/be/st/ss/hstgrades9through12.asp 

Conflict in Europe. Formation of the Third Reich. In 1933, German president Paul von Hindenberg named Adolf Hitler chancellor. As Civil Liberties began being limited, and the Nazification of Germany began in earnest, the Weimar Republic collapsed and the "Third Reich" began (in German, "Großdeutsches Reich"). Hitler had outlined his aims years earlier, in his book "Mein Kampf"(My Struggle). Hitler claimed that Germany had lost its powerful economy, morality, raw materials, land, and resources needed to help it develop as a nation. Hitler blamed "sub-human" peoples such as the Jews for his country's defeat. The superior German people needed "living room," and had a right to claim it. The book called for the elimination of the Jews, and the elimination of homosexuals, the mentally ill, and other "undesirable" elements of German society. Hitler also used this supposed German superiority, as well as German mistreatment by its victors after World War I, to justify the termination of the Treaty of Versailles. Military Buildup. Hitler began a buildup of the German military. In 1936, he tested German might by supporting the Fascists and German interests during the Spanish Civil War. Then Hitler and Benito Mussolini, the Fascist Dictator of Italy, created a coalition with the dictatorship which had come to power in Japan. The coalition of these three nations later came to be called the Axis. Appeasement policy. Many of the borders drawn up in the wake of Versailles were fragile. There were German nationalists in many other countries, including the post-war nation of Czechs and Slavs known as Czechoslovakia. In 1938, Hitler used alleged mistreatment of a German minority in another German-speaking nation to "annex" or take over Austria. Other nations were reluctant to interfere because of Hitler's claim that the relation between Germany and Austria was an internal German concern which had nothing to do with the rest of Europe. Then Hitler took control of a section of Czechoslovakia partially populated by Germans. This time, British Prime Minister Neville Chamberlain did interfere. In the wee hours of September 30, 1938, Chamberlain and Hitler signed an agreement ensuring that Germany would keep the territory it called "the Sudetenland", but would go no further in the country, nor take any German-populated areas in other European nations. The policy which sought to prevent another World War at almost any cost, including the cost of allowing a tyrant to gain more power, became known as "appeasement". Chamberlain called "the Munich Agreement" "Peace in our time." Hitler had no intention of keeping his word. In 1939, he took over the remainder of Czechoslovakia and demanded Poland. Great Britain and France agreed to come to Poland's aid. Then Germany signed another agreement with the Soviet Union, the confederation of Russia and its supporters under a Communist government. "The Nazi-Soviet Pact" was an agreement that the two nations would not fight each other. Both countries agreed to take parts of Poland, the Soviet Union securing the Baltic Sea port cities. (It had been a long-term interest of the Soviets to gain ice-free ports for winter trade.) Yet in private, Hitler was already planning to take the Soviet Union over in its turn. The Beginning of the War. Blitzkrieg. On the first day of September, 1939, Germany declared war on Poland. The British and French responded by declaring war on Germany two days later. The Germans used the tactic of "Blitzkrieg" ("lightning war") in Poland to defeat the Polish Army in as little as sixteen days as the British and French sat back in fear of a new World War. With little warning and no provocation, the German Air Force strafed Poland with top-speed planes and sent in tanks and racing artillery on the ground. The Polish had had no time to build the deep trenches seen in the First World War; the Germans had no need to use mustard gas. By the end of the first week of October, the Germans had gained control of half of Poland. The Soviets invaded from the East. With no time to defend themselves, the last Polish troops surrendered in early October. In the spring of 1940, Hitler attacked the nations of Denmark and Norway. Denmark surrendered, but British and French troops came to Norway's aid. Germany entered Belgium and the Netherlands on May Tenth, 1940. The Netherlands surrendered on May 15, though the province of Zeeland held out until the 18th. Belgium was overcome on May 28. On the same day, France recalled its troops from Norway, leaving Norway's fate to Germany. On June Fifth the Germans began their attack on France. To make matters worse, Mussolini declared war on France and Britain on June 10. The French government, meanwhile was taken over by a new Premier, who signed an armistice with Germany on June 17. Germany gained control of the northern part of France, and the Vichy French Government (so called because of the new French capital at Vichy) retained the south. The Italians had a small zone of occupation near the Franco-Italian border. Battle of Britain. Hitler's Germany was the supreme power on Continental Europe. Only the United Kingdom offered resistance. The Germans intended to invade the United Kingdom, but they first had to contend with the British Royal Air Force. The German Luftwaffe (Air Force) commenced the Battle of Britain in 1940. However, the British used the new technology of radar (Radio Detection and Ranging) to combat the Germans. In September, 1940, the Germans ended the Battle of Britain by indefinitely delaying all plans for invasion. Nonetheless, German airplanes continued to bomb several British cities until the middle of the next year. Invading the Balkans. Hitler expanded the Axis in the winter of 1940-1941 with the additions of Hungary, Romania, and Bulgaria. In April, 1941, Germany and Italy then attacked Yugoslavia, which surrendered within one week of invasion. Then Hitler and Mussolini turned to Greece, which collapsed by the end of April. By the end of 1942, most of Europe was under control of the Nazis or the Italians. Lend-Lease Act. In early 1941, the United States abandoned its neutrality and began to aid the British. "The Lend-Lease Act" allowed the President to lend or lease weapons worth over seven billion dollars to other nations. The first two years of the war overseas saw the American public broadly divided on the issue of potential involvement. Though the danger posed by Germany and Japan was generally recognized, millions of Americans felt that a strong, armed neutrality and oceanic defense without entering the war was the safest course. In contrast, President Roosevelt made it quite clear to those around him that he felt the United States would have to intervene on the Allied side, and planned and acted accordingly, initiating a war industrial buildup and proposing that the US become the "Arsenal of Democracy," supplying ammunition to Great Britain and its Allies. Conflict in the Pacific. Growing Tensions. On June 22, 1941, the Germans invaded the Soviet Union. The Pact between the two nations was dissolved, and the latter joined the Western Allies. Americans were very reluctant to start any conflict with Germany. Even in the fall of 1941, when shooting took place in the Atlantic between German U-boats and US ships, Roosevelt avoided escalation. After this, however, momentous events in the Pacific plunged the Americans into the war. The Great Depression had affected Japan as much as it had the Western powers. In 1931 a group of Japanese nationalists had assassinated the current Prime Minister of Japan, leading to a military dictatorship. In the same year Japan invaded the Chinese province of Manchuria. Although the government talked of nurturing international friendship, this was a take-over of Chinese national resources. In 1940 the Japanese marched into Indochina (present-day Southeast Asia), which had formerly belonged to the Dutch and to Vichy France. They now commanded the plantations responsible for the world's supply of rubber. The United States retaliated by attempting to prevent Japanese purchases of oil and steel. Tensions between Japan and the United States grew. A Date which will live in infamy. To secure resources and sea lanes for the Japanese islands, the Empire of Japan desired to neutralize the American Pacific Fleet, which was stationed at Pearl Harbor in Hawaii. On December Seven, 1941, the Japanese Air Force bombed the large American naval base, destroying or severely damaging over nineteen ships and 292 aircraft. The US Naval aircraft carriers, huge ships serving as mobile bases for airplanes, were then at sea and survived the attack. Yet its results were still dire: 2,403 American soldiers, sailors, and civilians were killed in the unprovoked strike. Japan made simultaneous strikes on Guam, Midway, and British bases. The next day, the United States Congress declared war on Japan, prompting Germany and Italy to in turn declare war on the United States. Japan continued with its Pacific operations by taking the American colonies of the Philippines, Guam, and Wake Island; the British colonies of Burma, Singapore, Malaya, and Borneo; and the Dutch colonies in the East Indies. Battle of Midway. An emboldened Japanese navy blundered in June of 1942, attacking "Midway Island" in the Pacific. After several days of aerial attacks on naval ships, American carrier-based planes defeated the Japanese ships so badly that their navy never recovered. Starving and weakened by disease, they held on for another month before surrendering. The Home Front. With the mass media of motion pictures and the radio, the American government was able to motivate and move individuals and groups more efficiently than ever before. The government had begun to draft eligible young men in 1940, and the draft was ramped up after Pearl Harbor. There was scattered opposition to this overwhelmingly popular measure among religious pacifists and the Nation of Islam, and among the American Communists until Germany attacked the Soviet Union. Under the streamlined system put in process during 1942, males eighteen years of age and older registered with the government. If the individual's name was drawn during an intermittent lottery, he received a postcard telling him to report to his local draft board, organized by the national Subscription Service. (In addition, many men, and some women, signed up without receiving a notice.) At the draft board men walked through the rooms of various doctors, psychiatrists, and military and civil authorities. An applicant could be refused if he was ill, even with such minor ailments as nearsightedness or flat feet. He might also apply for an exemption if he was the head of a family or in a steel or arms factory. There were also exemptions for pacifism, though the man had to have a letter from his priest or minister saying that he had a religious basis for his pacifism. Quakers and Seventh Day Adventists were both denominations known for non-violence. Conscientious objectors could become medics or work in some other peaceful operation in support of the war. Men who had no letter, and who were unwilling to support the War in peaceful ways, were often put in prison and derided by the public. From the draft board men went straight into Basic Training, from there to become soldiers, sailors or Marines, or airplane pilots. Women could become nurses or join various auxiliary forces such as the Woman's Army Corps (WACs). These were not permitted to fight, but in their capacity as communication and supervisory aides were often on the battlefield. With many men being shipped overseas to fight in the war, many employment opportunities opened up for women. Thousands of women began to support the war efforts by getting employed in many different factories producing tanks and planes and other sorts of weaponry. Some factories which had been devoted to peacetime production of cars and bicycles were converted to armament factories, and the nascent television stations were shut down. Sales of these items were halted "for the duration." The start of World War II brought the end of the Great Depression due to the supplies and men it takes to win a war. Part of the war effort was the incorporation of American civilians into efforts showing that they could also make a difference. Victory Gardens were even more abundant than in World War I. These supplemented food and gasoline rations, initiated with the aim of getting supplies to the troops. Public work sites, and even some schools and Boy and Girl Scouts, vied to round up paper, scrap metal, and other supplies. These scrap drives brought in a limited amount of useful material, but helped involve some people who would have felt useless in the big war drive. Some public citizens who had not been permitted to join the military were delegated as air raid wardens, making sure that homes were darkened at night. (Especially during the first few years of the war, there was a real fear that some enemy Axis plane might bomb the American mainland.) Hollywood, still reeling from the Depression, brought itself back with instructional films for the troops and propaganda for moviegoers at home. Japanese-Americans during WWII. In February 1942, the War Relocation Authority began to establish centers where Japanese-Americans, including those born in the United States and other citizens, were interned. Though this racial discrimination violated constitutional due process requirements, the Supreme Court ruled that such internment was lawful in 1944, when it decided "Korematsu v. United States". One potential factor in the decision to intern Japanese Americans was the results of the Niihau incident, where a few Japanese-Americans living on a Hawaiian island aided a downed Japanese pilot. Despite this, the vast majority of Japanese-American citizens, many of whom had never left America and grew up in America, were loyal to America. Additionally relatively few Japanese-Americans on Hawaii actually endured internment due to protests of local leaders on Hawaii that interment would cripple the economy of the islands. Some Japanese Americans were allowed to leave the camps to study or work. The 442nd Infantry Regiment was composed almost entirely of second generation Japanese-Americans with family in the internment camps, yet they served with distinction in the European Theater, becoming one of the most decorated units in American history relative to it's size and service length. Their actions were responsible for saving the Lost Battalion, American troops trapped behind German lines. Turning back the European Axis. During the summer and fall of 1941, the Germans kept up their amazing pace into the heart of Russia. By December they had reached Moscow, and Leningrad was under siege. The Soviets sent in reserve troops from Siberia, and launched a counter attack. It succeeded, and Moscow was saved. In the spring of 1942, Hitler ordered an attack into the Caucus Mountains, and Stalingrad. As they had done before, the Germans quickly advanced, breaking through the Russian lines. In Stalingrad, there was street to street, and house to house fighting. The Germans controlled over 90% of the city, but the Russians refused to surrender. A Russian reserve division encircled the Germans into the city, and 250,000 German soldiers were captured. It was one of the bloodiest battles in history. In 1943, the President of the United States for an unprecedented third term, Franklin D. Roosevelt, and the Prime Minister of the United Kingdom, Winston Churchill held a Conference at Casablanca. The two nations then set up a plan of action for the next stages of the war. Meanwhile, the Russians continued to hold back the Germans, inflicting a crucial and massive defeat on Hitler's armies at the battle of Stalingrad in the winter of 1942-43. After a further major Russian victory at Kursk the following summer, the Germans were forced into retreat back towards Europe. In Africa, Axis troops led by Erwin Rommel had pushed into Egypt, just 70 miles west of Alexandria. However, British troops led by General Montgomery decisively defeated the Italian and German troops at the Battle of El Alamein. They were pushed out of Egypt, all the way across Libya, and into Tunisia. In November 1942, the Americans launched operation Torch and drove the French troops out of Algeria and Morrocco. After a long battle with Axis troops in Tunisia, they were driven out of Africa in May 1943. The Allies then decided to invade Sicily, in hope of knocking Italy out of the war. In early July the invasion began. For the next month, the British and Americans led a bloody campaign in which Sicily was finally taken in early August. During the invasion Mussolini was overthrown and arrested. Hitler had him rescued and put him in charge of the new Italian Social Republic. Following the Invasion of mainland Italy in early September, the Italian government signed an armistice with the Allies. The fall of Italy signaled the beginning of the end of World War II. However, Mussolini was rescued by the Germans and had established an Italian Social Republic. Near the end of the War, Germany tries to fight for a last stand with the Allies. It became to be Germany's only hope for turning the war around. The battle took place in a 60 mile deep 40 mile wide "Bulge". Therefore giving the name of the battle, battle of the Bulge. After weeks of fighting in the cold winters the Allied forces came out victorious. Months after this battle the Allied forces had the Germans pushed all the way back into Berlin. Antisemitism and The Holocaust. The prejudice of racism in America -- though the term "racist" is a misnomer: we are all members of the human race -- was evident from the days of Christopher Columbus onward. Antisemitism was a powerful motivating force in American history. The limits on immigration set in the early 1920s was in part a reaction against [Jewish] immigrants from Eastern Europe. The prosecutors of "Red Summer" had a dread of Blacks, Jews, and atheists, and sometimes of a nightmare conglomeration of all three. But some antisemitism was an alien import. The Protocols of the Elders of Zion (Russian: "Протоколы сионских мудрецов", or "Сионские протоколы", see also other titles) is an antisemitic and anti-Zionist plagiarism and literary hoax first published in 1903 in Russian, in Znamya; it alleges a Jewish and Masonic plot to achieve world domination. A translation, taken as the literal truth, was published and popularized by the American industrialist Henry Ford. In the 1930s, the American Bund attempted to increase Nazi German influence, and to amplify the antisemitic messages coming from Berlin. President Roosevelt treated the topic with care. Ethnographers and other scientific experts were drafted into the War effort, emphasizing that Americans came from every ethnic background. Jews were in every branch of the service, and rabbis were among the chaplains brought to aid them. But what came to be known as "The Final Solution" -- Hitler's plan to eliminate what "Mein Kampf" had called "the Jewish Question" -- was unknown to the American public. American newspapers had printed accounts of German oppression of its Jewish citizens, from Kristallnacht in 1938 onward. They had occasionally mentioned persecution of its Gypsies, Slavs, and other "non-Aryans," and the dreadful punishments meted out to those who opposed the Nazi regime. In January, 1945, Soviet troops liberated Auschwitz, the largest of the Nazi Concentration Camps. On April Eleventh, 1945, Allies liberated the Death Camp at Buchenwald, near Weimar, Germany. On the 12th, several journalists arrived, including Edward R. Murrow, one of the most lauded journalists of the time. He sent out a broadcast for American audiences on the Fifteenth describing what he had seen and heard. "There surged around me an evil-smelling stink, men and boys reached out to touch me. They were in rags and the remnants of uniforms. Death already had marked many of them, but they were smiling with their eyes. I looked out over the mass of men to the green fields beyond, where well-fed Germans were ploughing..." On May third, Americans first saw newsreels of the camp. (Newsreels, film of weekly news reports shown in movie theaters, were major sources of information in the days before TV became popularized.) There they heard that the recent death rate had been about two hundred a day. They saw people in the last stages of malnutrition, disease, and "constant hard work, beatings, and torture." The camera also showed corpses of men stacked like cords of wood, and the crematoria where the dead had been burnt. Now for the first time the majority of Americans could guess at the extent of The Holocaust, one of the most ghastly episodes in the modern history of mankind. In April of 1933, three months after Hitler took power, the Nazis issued a decree ordering the compulsory retirement of "non-Aryans" from the civil service. This is known as the spark of the Holocaust. Before Germany was defeated, there were some eleven million people that had been slaughtered in the name of Nazi racial purity. Although the Jews were the favored targets and are the victims we most hear about when talking about the Holocaust, they were not the only victims. There were also millions of Russians, Poles, gypsies and others that were also murdered. Although the deprivation of the Jews started in the years following 1933, the mass killings didn't begin until 1941. The effect of this knowledge was augmented by the Nuremburg trials of 1945-1946. There many German officials, and some Concentration Camp governors, were put to trial for committing these murders, called crimes against humanity. American antisemitism has continued since that time, but it is at least officially frowned upon. The American eugenics movement also suffered a setback from which it has not yet recovered. Operation Overlord. In November, 1943, Prime Minister Churchill and President Roosevelt held another Conference at Tehran. Joseph Stalin, who held the title of General Secretary of the Communist Party of the USSR, but was actually a Dictator of the Soviet Union, joined them there. The three leaders agreed to a plan codenamed Operation Overlord, under which an attack would be launched on the northern coast of France from the English Channel. In preparation for an invasion of France, Hitler cut off all support for the German armies remaining in the Soviet Union. Thus disabled, the German Army was forced to withdraw from Russia in the winter of 1943-1944. On June 6, 1944 ("D-Day,") in the early morning hours, American and British paratroopers were dropped into Normandy. Hours later, American, British, French and Canadian soldiers landed at Normandy on the north coast of France. The troops landed near Caen, but Hitler wrongly felt that they would attack at a location to the north of that city. The Allies took advantage of Hitler's miscalculation; by the end of the month, the Allies had over eight hundred thousand soldiers in Normandy. Meanwhile, Russian troops, which had been on the defensive, began their offensive on German-controlled territories. In the middle of July, the Soviets won their first major victory by taking the territory of Belorussia. At this time, concern began to grow in the West about Soviet domination replacing German in eastern Europe, especially in Poland. Despite these worries, Roosevelt felt that he had little influence in that area over Stalin, whose armies were bearing a huge brunt of the fight. By the end of July, the Allies expanded their base at Normandy by breaking out into the rest of France. Pushing through the nation, the Allies had gone far enough to liberate the city of Paris on August 25. On September 11, some Allied troops entered Germany, taking Antwerp, Belgium on the way. German resistance then hardened, however. British Field Marshall Montgomery attempt to "end the war by '44" with Operation Market Garden, a plan to liberate Holland and bypass the German border defenses, failed. The British and American armies would make little more progress for the rest of 1944. Meanwhile, Russian troops pushed toward Germany, defeating Germany's Axis partners on the way. In August, Romania surrendered, followed by Bulgaria and Finland in September. Yalta and German Surrender. German Counteroffensive. Allied air bombing of German industries and cities had been ongoing and savage since 1943, but did not have the intended effect of crushing the German will to fight. Indeed, Hitler was able to field new advanced weapons in 1943-45, such as the world's first jet fighter aircraft, the V-1 flying bomb, the V-2 ballistic missile, and new types of tanks and submarines. The new weapons, however, proved of little use against Allied numbers and economic superiority, with American industrial production for the war effort massive and untouched by Axis attack. Germany forced millions of prisoners into slave labor, under the most brutal conditions, to keep its own war effort going. In December 1944, Germany launched a massive counter-attack on the light defended American positions in Belgium. The Germans hoped to cut off the Allied supply lines, however, after reinforcements arrived, the "Bulge"(today it is know as the Battle of the Bulge) was flattened out. Meanwhile, the Soviets were on the verge of entering Germany from the east en masse, having taken control of Poland. Hitler's troops were exhausted, millions dead or captured, and with the fall of the Romanian oil fields, German armies were running out of gasoline. A final call-up began of old men and boys for a last-ditch defense of Germany. Many German civilians fled, fearing the revenge the Russians would put on them after what the Germans had done in Russia. Thousands of German noncombatants were raped, and many of these were then killed. The Yalta Conference. In early February 1945, United States President Franklin D. Roosevelt, British Prime Minister Winston Churchill, and Soviet leader Josef Stalin, the three leaders of the anti-Nazi alliance, met at Yalta in the U.S.S.R. The leaders of the Allied powers, "The Big Three," were planning for the end of the war. This was a follow-up to the meeting of the three powers in November 1943 in the Tehran Conference. At Yalta the countries contended on what to do with Germany. Churchill and Britain wanted to protect their colonial possessions and to keep the Soviet Union from having too much power. Stalin and the Soviet Union wanted Germany to pay them to help start the rebuilding of their country. Stalingrad, a byword for Soviet military resistance of Hitler, was now in ruins. Two wars of German aggression had been too much, and the country needed to be permanently restrained. The United States wanted to influence Germany toward democracy and to keep the peace. The Yalta conference decided the division of defeated Germany into zones for reconstruction. The leaders agreed to punish Nazis for war crimes, including the Holocaust. Other topics included Soviet Russia's entry into the war against Japan, composition of the post-war government of Germany, voting arrangements in the new United Nations organizations, and the future of the liberated governments of Eastern Europe. The Yalta Declaration on Liberated Europe, which called for free elections and constitutional liberties in Eastern Europe, was the most controversial topic of the Yalta Conference. Race to Berlin. The Allies first attempted to reach the Rhine River in their quest to take over Germany. In March, this goal accomplished, the Americans and British opposed the Soviets in the "Race for Berlin". The Race determined who would control Berlin, a city that would prove important in the reconstruction of Germany. The Americans allowed the Soviets to win the Race for Berlin. Fierce fighting erupted in and around the city as motley German units made their last stand against the powerful army groups of Russian marshals Zhukov and Koniev. His capital surrounded and his loyal minions deserting him, Adolf Hitler committed suicide in his Berlin command bunker on April 30, 1945. Benito Mussolini was executed by Italian Partisans on April 28. The new leader of Germany, Karl Doenitz, agreed to surrender. On May 8, Germany formally signed an unconditional surrender, dissolving the Axis and leaving only Japan to be defeated. The End of the FDR Era. The public was never informed of the extent of President Roosevelt's paralysis. Though he admitted that he had suffered polio, no newsreels showed him pulling himself along with his crutches. At press conferences the photographers were led in when he was already in a chair or in his car seat. Compensation for the leg disability, and the physical activity the Roosevelts were known for, had made his upper arms strong and muscular. However, his physical strain combined with the burden of the Presidency to weary him. Speaking to Congress in the wake of Yalta, he made a rare admission of how much this affected him: "I hope that you will pardon me for this unusual posture of sitting down during the presentation of what I want to say, but . . . it makes it a lot easier for me not to have to carry about ten pounds of steel around on the bottom of my legs; and also . . . I have just completed a fourteen-thousand-mile trip." Roosevelt was ill, and had been frail even before his re-election in 1944. The Republican party had objected to the President even running for a fourth term. The Republican Candidate, Thomas E. Dewey, accused him of running a corrupt autocracy and of feigning good health. In response to accusations of corruption, Roosevelt dropped former Vice President Henry A. Wallace and picked up newcomer Harry S. Truman, who had risen in a War anti-fraud committee in the Senate. He also campaigned vigorously in the cities, which may have hastened his death. Roosevelt died of a stroke on April 12th, 1945. After FDR's death Harry Truman became president. The day after Roosevelt's death Truman sought out old friends to ask for their help in this "terrible job." The Atomic Bomb and the End of World War II. Island Hopping &amp; Kamikaze. After the Battle of Midway the United States surely retook the Asian Pacific nations, fighting the Japanese empire island by island, by gun, by shell, and by flame thrower. Though the Japanese continued to fight, its armed forces were in a hopeless situation. Its armaments became fewer. Kamikaze (the Japanese word for "Sacred Wind"), pilots who hoped to destroy American ships by intentionally crashing their own planes into them, became more dominant in Japanese opposition. The Manhattan Project. The United States had become involved in a competition with Nazi Germany to find technologies which would win the war. Among these technologies was nuclear fission, the splitting of a highly electronegative atom into smaller ones, which gives off ten times as much power. With cooperative or captive scientists and slave labor, the Nazis were attempting to use uranium to create an atomic bomb. During the same time, America had become the home for dissident German and Jewish scientists. Native-born and newly-emigrated American physicists were also talking about the possibilities of fission. One of the latter, Albert Einstein, sent a letter to President Roosevelt in 1939 explaining the developments in the nuclear chain reactions which would result in an atomic explosion. In 1942, a number of the top minds in physics were relocated to a secret location in New Mexico. These men volunteered to work at what was code-named the Manhattan Project, a secret attempt to bring to fruition what Einstein had posited in 1939. At its height the Manhattan Project employed more than 600,000 workers, the majority of them not realizing what they were really working for. After spending more than two billion dollars, the Project created the first atomic bomb.On July 16, 1945, this bomb was successfully tested in New Mexico. In the interim, Germany had been defeated. President Roosevelt had died, and had been replaced in office by his Vice President, Harry S. Truman. President Truman listened to experts who forecast two more years of war against Japan, and the prospect of a bloody American invasion. Truman chose to use the atomic bomb instead of an invasion. Hiroshima and Nagasaki. On August sixth, an A-bomb with the name Little Boy was dropped on Hiroshima, Japan, by a B-29 aircraft piloted by Col. Paul Tibbets. The explosion formed a mushroom cloud. Dust and debris shot into the sky and was seen miles away. Many people died instantly, but others fell sick a few days later from effects on the radiation on living tissue. The American press was told about the detonation of the bomb on a population of civilians, but not about its full effects. The Japanese government did not give in. On August ninth a bomb called Fat Man was dropped on the weapon-producing city of Nagasaki, again upon civilians. Together, the bombs killed over one hundred thousand people. Between the two bombings, meanwhile, the Soviet Union had joined in the war on Japan. The Americans threatened a third bombing on Tokyo, though they had not yet had time to create a new bomb. Now Japan unconditionally surrendered, officially ending World War II. Officials signed a treaty for surrender on September 2, 1945, aboard the USS Missouri. Death Toll. The death toll of the Second World War was greater than that of the First. At least sixty-one million people died from the Allied nations of the Soviet Union; the United States and its colonies; Great Britain, its colonies and Canada; France and its colonies; the Netherlands and its colonies; Belgium and its colonies; and Poland, Norway, and Greece. In contrast, the main Axis powers of Germany, Italy and Japan only suffered twelve million casualties. 

Eisenhower. Civil Rights Movement under Eisenhower and Desegregation. The first events that would spark off the entire Civil Rights movement happened during the Eisenhower administration. In the south, there were many statewide laws that segregated many public facilities ranging from buses to water fountains. Southern African Americans now felt that their time had come to enjoy American democracy and they fought hard to end southern segregation policies. "Brown v. Board of Education". In 1952, seven year old Linda Brown, of Topeka, Kansas, wasn't permitted to attend a white-only elementary school that was only a few blocks from her house. In order to attend her coloreds-only school, Brown had to cross dangerous railroad tracks and take a bus for many miles. Her family sued the Topeka school board and lost, but appealed the case all the way to the Supreme Court. "Brown v. Board of Education of Topeka, Kansas" came to the Supreme Court in December 1952. In his arguments, head lawyer for the NAACP, Thurgood Marshall, challenged the "Separate But Equal" doctrine established in "Plessy v. Ferguson" in 1896. He argued that schools could be separate, but never equal. On May 17, 1954, the Court gave its opinion. It ruled that it was unconstitutional to segregate schools, and ordered that schools integrate "with all deliberate speed." Central High Confrontation. Integration would not be easy. Many school districts accepted the order without argument, but some, like the district of Little Rock, Arkansas, did not. On September 2, 1957, the day before the start of the school term the Arkansas Governor, Orval Faubus, instructed the National Guard to stop any black students entering the school. He claimed this was to protect the property against violence planned by integration protesters. The federal authorities intervened and an injunction was granted preventing the National Guard from blocking the school and they were withdrawn on September 20. School restarted on 23 September, with the building surrounded by local police officers and nearly one thousand protesters. The police escorted nine black students, later known as the Little Rock Nine, into the school via a side door. When the crowd discovered the students had entered the building, they tried to storm the school and the black students were hurried out around lunch time. Congressman Brooks Hays and the Little Rock mayor, Woodrow Mann, asked the federal government for more help. On September 24, Mann sent a message to President Eisenhower requesting troops. Eisenhower responded immediately and the 101st Airborne Division was sent to Arkansas. In addition, the President brought the Arkansas National Guard under federal control to prevent its further use by the Governor. On September 25, 1957, the nine black students finally began their education properly, protected by 1,000 paratroopers. Montgomery Bus Boycott. On December 1, 1955, Rosa Parks, a seamstress and secretary of the Montgomery, Alabama, the chapter of the NAACP, boarded a city bus with the intention of going home. She sat in the first row of seats in the "colored" section of the segregated bus. At the next stop, whites were among the passengers waiting to board but all seats in the "white" front dividing the black and white sections to accommodate the racial makeup of the passengers at any given moment. So he ordered the four blacks sitting in the first row of seats in the "colored" section to stand and move to the rear of the bus so the waiting whites could have those seats. Three of the passengers complied; Mrs. Parks did not. Warned again by the driver, she still refused to move, at which point the driver exited the bus and located a policeman, who came onto the bus, arrested Mrs. Parks, and took her to the city jail. She was booked for violating the segregation ordinance and was shortly released on bail posted by E. D. Nixon, the leading local civil rights activist. She was scheduled to appear in municipal court on December 5, 1955. Mistreatment of African Americans on Montgomery's segregated buses was not uncommon, and several other women had been arrested in similar situations in the months preceding Parks's. However, Mrs. Parks was especially well-known and well-respected within the black community, and her arrest particularly angered the African Americans of Montgomery. In protest, community leaders quickly organized a one-day boycott of the buses to coincide with her December 5 court date. An organization, the Montgomery Improvement Association, was also created, and the new minister of the Dexter Avenue Baptist Church, the 26-year-old Martin Luther King, Jr., was selected as the MIA's president. Word of the boycott spread effectively through the city over the weekend of December 3–4, aided by mimeographed fliers prepared the Women's Political Council, by announcements in black churches that Sunday morning, and by an article in the local newspaper about the pending boycott, which had been "leaked" to a reporter by E. D. Nixon. On the morning of Mrs. Parks's trial, King, Nixon, and other leaders were pleasantly surprised to see that the boycott was almost 100 percent effective among blacks. And since African Americans made up 75% of Montgomery's bus riders, the impact was significant. In city court, Mrs. Parks was convicted and was fined $10. Her attorney, the 24-year-old Fred D. Gray, announced an appeal. That night, more than 5,000 blacks crowded into and around the Holt Street Baptist Church for a "mass meeting" to discuss the situation. For most in the church (and listening outside over loudspeakers), it was their first time to hear the oratory of Martin Luther King, Jr. He asked the crowd if they wanted to continue the boycott indefinitely, and the answer was a resounding yes. For the next 381 days, African Americans boycotted the buses, while the loss of their fares drove the Chicago-owned bus company into deeper and deeper losses. However, segregationist city officials prohibited the bus company from altering its seating policies, and negotiations between black leaders and city officials went nowhere. With bikes, carpools, and hitchhiking, African Americans were able to minimize the impact of the boycott on their daily lives. Meanwhile, whites in Montgomery responded with continued intransigence and rising anger. Several black churches and the homes of local leaders and ministers, including those of Nixon and King, were bombed, and there were numerous assaults by white thugs on African Americans. Some 88 local black leaders were also arrested for violating an old anti-boycott law. Faced with the lack of success of negotiations, attorney Gray soon filed a separate lawsuit in federal court challenging the constitutionality of the segregated seating laws. The case was assigned to and testimony was heard by a three-judge panel, and the young Frank M. Johnson, Jr., newly appointed to the federal bench by Republican President Dwight D. Eisenhower, was given the responsibility for writing the opinion in the case. Johnson essentially ruled that in light of the 1954 Brown v. Board of Education decision by the U.S. Supreme Court, there was no way to justify legally the segregation policies, and the district court ruling overturned the local segregation ordinance under which Mrs. Parks and others had been arrested. The city appealed, but the U.S. Supreme Court affirmed the lower court ruling, and in December 1956, city officials had no choice but to comply. The year-long boycott thus came to an end. The Montgomery Bus Boycott made Mrs. Parks famous and it launched the civil rights careers of King and his friend and fellow local minister, Ralph Abernathy. The successful boycott is regarded by many historians as the effective beginning of the twentieth-century civil rights movement in the U.S. Foreign Policy. In addition to his desire to halt the advance of “creeping socialism” in U.S. domestic policy, Eisenhower also wanted to “roll back” the advances of Communism abroad. After taking office in 1953, he devised a new foreign policy tactic to contain the Soviet Union and even win back territory that had already been lost. Devised primarily by Secretary of State John Foster Dulles, this so-called New Look at foreign policy proposed the use of nuclear weapons and new technology rather than ground troops and conventional bombs, all in an effort to threaten “massive retaliation” against the USSR for Communist advances abroad. In addition to intimidating the Soviet Union, this emphasis on new and cheaper weapons would also drastically reduce military spending, which had escalated rapidly during the Truman years. As a result, Eisenhower managed to stabilize defense spending, keeping it at roughly half the congressional budget during most of his eight years in office. The doctrine of massive retaliation proved to be dangerously flawed, however, because it effectively left Eisenhower without any options other than nuclear war to combat Soviet aggression. This dilemma surfaced in 1956, for instance, when the Soviet Union brutally crushed a popular democratic uprising in Hungary. Despite Hungary’s request for American recognition and military assistance, Eisenhower’s hands were tied because he knew that the USSR would stop at nothing to maintain control of Eastern Europe. He could not risk turning the Cold War into a nuclear war over the interests of a small nation such as Hungary. The Warsaw Pact and NATO. 1955 saw the division of Europe into two rival camps. The westernized countries of the free world had signed NATO 1949 and the eastern European countries signed the Warsaw pact. NATO. [[File:Truman signing North Atlantic Treaty.jpg|thumb|President Truman signing the North Atlantic Treaty.]] The North Atlantic Treaty Organization was created as a response to the crisis in Berlin. The United States, Britain, Canada, France, Portugal, Italy, Belgium, Luxembourg, Norway, Denmark, Iceland, and the Netherlands founded NATO in April 1949, and Greece, Turkey and West Germany had joined by 1955. The countries agreed that "an armed attack against one or more of [the member states] in Europe or North America shall be considered an attack against them all," and was created so that if the Soviet Union eventually did invade Europe, the invaded countries would have the most powerful army in the world (the United States' Army) come to their defense. When the Korean War broke out, NATO drastically raised its threat level because of the idea that all the communist countries were working together. As the number of communist countries grew and grew, so did the NATO forces. Greece and Turkey eventually joined NATO in 1952. The USSR eventually decided to join NATO so that there would be peace, but NATO declined them because they thought that the USSR would try to weaken them from the inside. The Warsaw Pact. The Soviet Union responded in to the addition of West Germany to NATO 1955 with its own set of treaties, which were collectively known as the Warsaw Pact. Warsaw Pact was also known as “The Treaty of Friendship”. The Warsaw Pact allowed East Germany, Poland, Czechoslovakia, Hungary, Albania, Romania, and Bulgaria to function in the same way as the NATO countries did. The Soviet Union used this Warsaw Pact to combine the military forces unified under it. The Pact was supposed to make all the countries in it, equal. However, the Soviet Union took a little advantage of this by using the allied countries military wherever they wanted. Unlike NATO, Warsaw forces were used occasionally. CIA. As an alternative, Eisenhower employed the CIA to tackle the specter of Communism in developing countries outside the Soviet Union’s immediate sphere of influence. Newly appointed CIA director Allen Dulles (the secretary of state’s brother) took enormous liberties in conducting a variety of covert operations. Thousands of CIA operatives were assigned to Africa, Asia, Latin America, and the Middle East and attempted to launch coups, assassinate heads of state, arm anti-Communist revolutionaries, spread propaganda, and support despotic pro-American regimes. Eisenhower began to favor using the CIA instead of the military because covert operations didn’t attract as much attention and cost much less money. [[File:Operationajax.jpg|thumb|The 1953 coup in Iran]] A CIA-sponsored coup in Iran in 1953, however, did attract attention and heavy criticism from liberals both at home and in the international community. Eisenhower and the Dulles brothers authorized the coup in Iran when the Iranian government seized control of the British-owned Anglo-Iranian Oil Company. Afraid that the popular, nationalist, Soviet-friendly prime minister of Iran, Mohammed Mossadegh, would then cut off oil exports to the United States, CIA operatives convinced military leaders to overthrow Mossadegh and restore Mohammed Reza Shah Pahlavi as head of state in 1953. Pahlavi returned control of Anglo-Iranian Oil to the British and then signed agreements to supply the United States with almost half of all the oil drilled in Iran. The following year, a similar coup in Guatemala over agricultural land rights also drew international criticism and severely damaged U.S.–Latin American relations. Vietnam. [[File:Ho Chi Minh (third from left standing) and the OSS in 1945.jpg|thumb|Hồ Chí Minh (Standing, third from left), posing with members of the Office of Strategic Services in 1945, a predecessor organization to the CIA. Ho Chi Minh had worked with Americans during World War II to collect intelligence against the Japanese, and hoped America could help create a Vietnam independent from French rule.]] In 1945 many colonies, including French Indochina, hoped for independence following the War. When Japan surrendered to the allies a Vietnamese man, Ho Chi Minh, declared independence for Vietnam in Hanoi, quoting America's own Declaration of Independence in the very first lines of his speech in hopes of gaining American support for a Vietnam free from French rule. However, with containing communism in Europe being seen as a more important issue, America took a stance of neutrality from 1946 to 1950. In the early 50's, Vietnam was rebelling against French rule. America saw Vietnam as a potential source of trouble, as rebels (known as the Việt Minh) led by Communist leader Ho Chi Minh were gaining strength. America loaned France billions of dollars to aid in the war against the Vietnamese rebels, but despite the aid, France found itself on the verge of defeat, and appealed to America for troops, but America refused, fearing entanglement in another costly Korean War, or even a war with all of communist Asia. France surrendered, and the VietMinh and France met in Geneva, Switzerland to negotiate a treaty. Vietnam was divided into two countries: the Vietminh in control of the North and the French-friendly Vietnamese in control of the South. In 1956, the two countries would be reunited with free elections. Eisenhower worried about South Vietnam. He believed that if it also fell to the Communists, many other Southeast Asian countries would follow, in what he called the "domino theory". He aided the Southern government and set up the Southeast Asia Treaty Organization (SEATO) in 1954. The nations included in the alliance were America, Great Britain, France, Australia, Pakistan, the Philippines, New Zealand and Thailand, and they all pledged to fight against "Communist aggressors". Cuban Revolution. [[File:CheyFidel.jpg|thumb|upright|Che Guevara and Fidel Castro.]] In 1958 and 1959, anti-American feeling became a part of the growing Cuban revolution. In January 1959, the dictator of Cuba, Fulgenicio Batista, was overthrown by the rebel leader Fidel Castro, who promptly became the leader of Cuba. At first, America supported Castro because of his promises of democratic and economic reforms. But relations between the two countries became strained when Cuba began seizing foreign-owned land (which was mostly U.S. owned) as a part of its reforms. Soon, Castro's government was a dictatorship, and was being backed by the Soviet Union. In 1961, Eisenhower cut diplomatic ties with Cuba, and relations with the island nation have been difficult ever since. Suez Canal. [[File:Port Said from air.jpg|thumb|British and French forces make a move on Port Said during the crisis.]] Back in 1948, Israel was created as a sanctuary of sorts for the displaced Jews of the Holocaust. At the same time, many Arabs living in the area were displaced. Tensions had been high in the Middle East ever since Israel had been attacked just after its founding. The stage was set for superpower involvement in 1956; the United States backed Israel, the Soviet Union backed the Arabs, and the Egyptian president Gamal Abdel Nasser had nationalized, or brought under Egypt's control, the Suez Canal, which had previously belonged to Britain. France and Britain worried that Egypt would decide cut off oil shipments between the oil-rich Middle East and western Europe, so that October they invaded Egypt, hoping to overthrow Nasser and seize the canal. Israel, upset by earlier attacks by Arab states, agreed to help in the invasion. U.S. and Soviet reactions to the invasions were almost immediate. The Soviets threatened to launch rocket attacks on British and French cities, and the United States sponsored a United Nations resolution for British and French withdrawal. Facing pressures from the two powers, the three invaders pulled out of Egypt. To ensure stability in the area, United Nations troops were sent to patrol the Egypt-Israel border. Space Race. [[File:President Dwight D. Eisenhower, Dr. von Braun and Others.jpg|thumb|President Eisenhower meets with Dr. Wernher von Braun in 1960.]] The Space Race has its origins in an arms race between the United States and the Soviet Union. After the Soviet Union tested its first atomic bomb in September 1949, the fear of nuclear war began to spread. The ensuing arms race led to the creation of intercontinental ballistic missiles (ICBMs), long-range rockets designed to deliver nuclear warheads from land- or submarine-based launch sites. The first successful ICBM test flight was on August 21, 1957, when the USSR launched an ICBM with a dummy payload over 4,000 miles to an isolated peninsula on Russia’s east coast that had been declared a military zone (Kamchatka Peninsula, which remained closed to civilians from 1945–1989). On October 4, 1957, the Soviets successfully put the first man-made satellite, "Sputnik", into orbit. Americans were horrified. They feared that the Soviets were using the satellite to spy on Americans, or even worse, that the Soviet Union might attack America with nuclear weapons from space. America responded with the launch of its own satellite, "Vanguard". Hundreds of spectators gathered, only to watch the satellite rise only a few feet off the launch pad, and then explode. The failure spurred the government to create a space agency, the National Aeronautics and Space Administration (NASA). NASA succeeded in launching the "Explorer" in 1958, and thus, the Space Race was initiated. With the creation of Project Mercury, a program to put an astronaut in space, America was pulling ahead. Nonetheless, the USSR was the first to put a man in space, when Yuri Gagarin was launched into orbit in 1961. For the next 14 years, the U.S and the Soviet Union would continue to compete in space. Many Americans were frightened during the start of the space race because this gave the Soviet Union a better ability to launch a surprise attack on the states. The United States of course immediately jumped on the bandwagon and did its best to become the first country to land on the moon. Domestic Policies. [[File:2008-0831-TheGreenbrier-North.jpg|thumb|Under the Eisenhower Administration [[w:Project Greek Island|Project Greek Island]] was started to create a covert bunker beneath a resort near Washington DC where the government could continue to operate in the event of Nuclear War. Completed in 1962, it's existence was revealed to the public in 1992 by a Washington Post expose.]] Alaska and Hawaii were admitted as states during the Eisenhower administration. Eisenhower established the Interstate Highway System in 1956 to bolster the economy, and to facilitate rapid mobilization of defense forces in the event of an emergency. At the end of his administration, Eisenhower warned of the dangers of a Military Industrial Complex, fearing undemocratic influence from a rapidly growing defense sector. Everyday life in the 1950's. [[File:DowneyMcdonalds.jpg|thumb|The third McDonald's restaurant, and the oldest operating McDonald's. The aesthetics of this location are nearly unchanged from when it opened in 1953.]] With the rise of television ownership came the rise of television shows. Emerging genres like Sitcoms (Situational comodies), talk shows, and game shows enjoyed widespread popularity. Easy to cook, ready made meals called "TV dinners" were introduced during this time The construction of the interstate system firmly established a car culture in America. Fast Food had existed prior to the 1950's, but the 1950's were when it really became a phenomenon, when franchising locations of popular restaurants became common, eventually leading to the formation of chains with uniform quality and items nationwide. In the 1950's Pizza became an common staple outside of the Italian-American community, spurred by American troops returning from Italy, references to pizza in popular culture, and the establishment of chain pizza restaurants. Rise of the Middle Class. In the late 1940s and 1950’s led to a rise of the middle class in the United States. With troops coming home from the war soldiers were quick to start families. The GI bill allowed for upward mobility for veterans. With free college and over four billion dollars given to trips the economy continued to succeed. With war savings and a host of new consumer goods on the market, America quickly turned into a consumer market. The best example of this would be automobiles and the television. In 1955, $65 billion was spent on automobiles. This represented 20% of the Gross National Product. In 1950 50% of American homes had a television. By 1960 this number was raised to 90%. Rock and Roll. [[File:Chuck Berry 1957.jpg|thumb|upright|Chuck Berry in 1957.]] The term “rock and roll” was originally a nautical phrase referring to the motion of a ship at sea. In the early 20th century, it gained a religious connotation (referring to the sense of rapture felt by worshippers) and was used in spirituals. After this, “rocking and rolling” increasingly became used as a metaphor for sex in blues and jazz songs. The origins of rock and roll lie primarily in electric blues from Chicago in the late 1940s, which was distinguished by amplification of the guitar, bass, and drums. Electric blues was played by artists like Muddy Waters, Willie Dixon, and Buddy Guy, who were recorded by Leonard and Phil Chess at Chess Records in Chicago. They inspired electric blues artists in Memphis like Howlin’ Wolf and B.B. King, who were recorded by a Memphis-based record producer named Sam Phillips, also the owner of Sun Records. He later discovered Elvis Presley in 1954, and he also recorded early songs by Jerry Lee Lewis, Johnny Cash, Roy Orbison, and Carl Perkins. Rhythm and blues artists such as Little Richard, Chuck Berry, Bo Diddley, and Ray Charles incorporated electric blues as well as gospel music. In the early 50s, R&amp;B was more commonly known by the blanket term “race music”, which was also used to describe other African-American music of the era such as jazz and blues. Billboard didn’t replace the “race records” category with “rhythm and blues” until 1958. Doo-wop was a mainstream style of R&amp;B, with arrangements favoring vocal harmonies. Other types of music that contributed to early rock and roll include African-American spirituals, also known as gospel music, and country/folk, which was primarily made by poor whites in the South. Arguably, the first rock and roll song ever made is “Rocket 88”, recorded at Sam Phillip’s studio in Memphis in March 1951. It was credited to “Jackie Brenston and his Delta Cats,” a band that didn’t actually exist—the song was put together by Ike Turner and his band, the Kings of Rhythm. Jackie Brenston was the vocalist on the song, who also played saxophone in the band. What really sets “Rocket 88” apart is the distorted guitar sound: it was one of the first examples of fuzz guitar ever recorded. The amplifier they used to record the song was damaged on the way from Mississippi to Memphis. They tried to hold the cone in place by stuffing the amplifier with newspaper, which created the distortion. Sam Phillips liked the sound and decided to keep it in the song. Although “Rocket 88” was recorded by Sam Phillips, it wasn’t released by Sun Records, which didn’t exist until 1952. From 1950-1952, Phillips ran the Memphis Recording Service, where he would let amateurs perform and then sell the recordings to large record labels. He sold “Rocket 88” to Chess Records, which released predominately blues, gospel, and R&amp;B. The Chess brothers started Checker Records in 1952, because radio stations would only play a certain number of tracks from each label. Alan Freed (also known as “Moondog”) was a radio DJ who started playing R&amp;B records on WJW in Cleveland in 1951. He is credited with introducing rock and roll to a wide audience for the first time, as well as being the first to use the phrase “rock and roll” as the name of the genre. He also promoted and helped organize the first major rock and roll concert, The Moondog Coronation Ball, which occurred on March 21, 1952. The concert was so successful that it became massively overcrowded – there was a near-riot and it had to be shut down early. Freed’s popularity soared, and he was immediately given more airtime by the radio station. His promotion of rock ‘n’ roll is one of the main reasons it became successful, and in recognition of his contributions to the genre, the Rock and Roll Hall of Fame was built in Cleveland. “Payola” refers to the practice of record company promoters paying radio DJs to play their recordings in order to boost their sales. Payola had been commonplace since the Vaudeville era in the 1920s, but it became a scandal in the 1950s due to a conflict between the American Society of Composers, Authors, and Publishers (ASCAP) and radio stations. Prior to 1940, ASCAP had made huge amounts of money from the sales of sheet music, but when radio started gaining popularity, recorded music became more profitable than sheet music. ASCAP demanded large royalty payments from radio stations that played their recordings. Instead, stations boycotted ASCAP recordings and created their own publishing company called Broadcast Music Incorporated (BMI). ASCAP tended to ignore music composed by black musicians or “hillbillies”, which gave BMI control of these areas. When rock ‘n’ roll became more and more popular, BMI became more and more successful. ASCAP (in addition to many others) believed that rock ‘n’ roll was the music of the devil, that it was brainwashing teenagers, and that it would never have been successful without payola. This was just after the quiz show scandal (when it was found that certain shows were rigged), and ASCAP urged the House Legislative Committee which had investigated that scandal to look into payola. The hearings that followed destroyed Alan Freed’s career, although it didn’t eliminate rock ‘n’ roll altogether as ASCAP had hoped. [[File:Elvis Presley promoting Jailhouse Rock.jpg|thumb|upright|Elvis Presley promoting Jailhouse Rock in 1957.]] Several factors contributed to the decline of early rock and roll. Chuck Berry and Jerry Lee Lewis were both prosecuted in scandals involving young women. Elvis Presley was inducted into the U.S. Army in 1958, and after training at Fort Hood, he joined the 3rd Armored Division in Germany, where he would remain until 1960. Three rock and roll musicians –- Buddy Holly, Ritchie Valens, and “The Big Bopper” –- died in a plane crash on February 3, 1959 (“The Day the Music Died”). Little Richard retired from secular music after a religious experience. He ran a ministry in Los Angeles, preached across the country, and recorded gospel music exclusively until 1962. Rock and roll music is associated with the emergence of a teen subculture among baby boomers. Teenagers bought records and were exposed to rock and roll via radio, jukeboxes, and television shows like American Bandstand, which featured teenagers dancing to popular music. It also affected movies, fashion trends, and language. The combination of white and black music in rock and roll –- at a time when racial tensions were high and the civil rights movement was in full-swing –- provoked strong reactions among the older generation, many of whom worried that rock and roll would contribute to social delinquency among teenagers. However, it actually encouraged racial cooperation and understanding to some extent—rock and roll was a combination of diverse styles of music made by different races, and it was enjoyed by both African-American and Caucasian teens. 

A topological space is a set X, and a collection of subsets of X, C such that both the empty set and X are contained in C and the union of any subcollection of sets in C and the intersection of any finite subcollection of sets in C are also contained within C. The sets in C are called open sets. Their complements relative to X are called closed sets. Given two topological spaces, X and Y, a map f from X to Y is continuous if for every open set U of Y, f−1(U) is an open set of X. 

Vectors are commonly used in physics and other fields to express quantities that cannot be accurately described by a scalar. Scalars are simply the value of something in a single dimension - a real number. For example, one might say that they have driven 5 kilometers, that an hour has elapsed, or that something's mass is 20 kilograms. In every one of these cases, there has been exactly one value stated. However, we might have more information we wish to give. Take the example of driving 5 kilometers. In this case, it may be useful to know how far you drove, but it might also be equally important which "direction" you drove, such as 5 kilometers due east. Now, given your starting point, exactly where you drove can be determined. Definitions. Vectors can be described mathematically by using trigonometry. We can define a vector to be an ordered pair consisting of a magnitude and a direction. In this diagram, "r" is the magnitude of this vector and θ is the direction. Notice, now, that we have moved horizontally "r" cos(θ) and vertically "r" sin(θ). These are called the "x-component" and the "y-component", respectively. We can also write a vector conveniently in terms of the x and y component. We write formula_1 for vectors. In some texts, you may see the vector written sideways, like ("x", "y"), but when you write it will help "greatly" to write them downwards in columns. In print we commonly use bold vectors, but since you probably don't have a pen that writes in bold print, underline your vectors, i.e. write v, or put a tilde underneath your vectors. Occasionally in Physics, you may see vectors written with an arrow pointing right. Notice that vectors need not have two components. We can have 2 or 3 or "n" or an infinite number of components. We write the set of all vectors with 2 real number components as R2; likewise for 3, "n", or infinite number of components. For components with complex numbers, we write C. Polynomials are "vectors" too - we'll look at notation for the set of polynomials later. For a reason why we do this, see Set theory for an explanation. Stretching and shrinking. We can define some actions on vectors. What will happen if we extend the vector? Or what will happen if we shrink the vector? The vector's "direction" doesn't change, only its length -- its magnitude. The action we perform to stretch or shrink a vector is that we multiply its magnitude by some amount. We refer to doing this as "scalar multiplication": we multiply the "vector" by a "scalar" real number. Scalar multiplication. For scalar multiplication, we simply multiply each component by the scalar. We commonly use Greek letters for scalars, and English letters for vectors. So for a scalar value of λ and a vector v defined by "r" and θ, the new vector is now λ"r" and θ. Notice how the direction does not change. Example. Say we have formula_2 and we wish to double the magnitude. So, formula_3. Addition of vectors. Simply, to add two vectors, you must add the respective x-components together to obtain the new x-component, and likewise add the two y-components together to obtain the new y-component. Example. Say we have formula_4 and we wish to add these. So, formula_5. Subtraction of vectors. The operation of subtraction on two vectors, a and b, a-b, can also be written as a+(-1)b. Therefore, we can use scalar multiplication to find the value of (-1)b, then use vector addition to find our solution. Complex numbers as vectors. ../Complex numbers/ can be represented in the form formula_6 or equivalently formula_7 , or in other words, a vector with magnitude formula_8 and direction formula_9. On the complex plane, this vector has a "real" x-component and an "imaginary" y-component. See ../Complex numbers/ for more information. Lines and planes. We can form the equations of lines and planes using vectors. Let's see how we can do this. Vector equation of the line. Consider a vector formula_10. Let's consider the following: If we have the equation λv, it is clear that for each choice of λ we choose, we get a different point on the line "y"=2"x". We can now generalize this idea into the "vector equation of the line" (and it is not restricted to 2 dimensions either). The vector equation of a line is given by where v is a vector parallel (which then, could lie) on the line. λ then, is the unknown in the equation. x is then the dependent vector variable. Vector equation of the plane. Now consider a plane. If we have two nonparallel vectors lying on the plane and we add them, we can add a linear combination (that is, add the two vectors, which are multiplied only by scalars) to choose some other vector. The set of all vectors under linear combinations of these two vectors form a plane. More simply, if we have two nonparallel vectors a and b we can form any other vector parallel to a and b by: where λ1 and λ2 are both scalars. Further algebra and geometry of vectors. There are other operations on vectors which we can perform. These operations we will consider have very real and significant geometric meanings. Magnitude. The "magnitude" of a vector is its length in R+ The dot product. The "dot product" of two vectors is defined as the sum of the products of the components. Symbolically we write For example, Properties of the dot product. If we have a and b as vectors, where "c" is a scalar. Geometry of the dot product. The dot product of two vectors has an alternate form: If we pick a vector c=a-b to form a triangle, we can show that these two forms are indeed equivalent by trigonometry. The angle θ then is important, as it shows that the dot product of two vectors is related to the angle between them. More specifically, we can calculate the dot product of two vectors - if the dot product is zero we can then say that the two vectors are perpendicular. For example, consider simply Plot these vectors on the plane and verify for yourself that these vectors are perpendicular. Cross product. The "cross product" is a more complicated product to define, but has a nice geometric property. We will only look at the cross product in three dimensions, since it is the most commonly used in three dimensions and it is difficult to define in greater dimensions. For a vector with three components, the cross product is defined as where If you have not done Matricies before, here is a formula to work out from above... = i formula_20 - jformula_21+ k formula_22 Properties of the cross product. The cross product has some properties which is easily verified from the above definition, and Geometric properties of the cross product. The cross product has some interesting geometric properties. If a and b are two vectors, a×b is the vector perpendicular to both. Now if we have two vectors, we have "two" choices of vector perpendicular to a and b - if we switch the order of the cross product we obtain the other vector. The magnitude of the cross product of two vectors is the area of the parallelogram formed by these two vectors. The "scalar triple product", a·(b×c) is the volume of the paralleliped formed by these three vectors. 

Remember, SN1 and SN2 are the same reaction just undergoing a different mechanism. 

Alcohols are the family of compounds that contain one or more hydroxyl (-OH) groups attached to a single bonded alkane. Alcohols are represented by the general formula -OH. Alcohols are important in organic chemistry because they can be converted to and from many other types of compounds. Reactions with alcohols fall into two different categories. Reactions can cleave the R-O bond or they can cleave the O-H bond. Ethanol (ethyl alcohol, or grain alcohol) is found in alcoholic beverages, CH3CH2OH. =Preparation= In the Alkenes section, we already covered a few methods for synthesizing alcohols. One is the hydroboration-oxidation of alkenes and the other is the oxymercuration-reduction of alkenes. But there are a great many other ways of creating alcohols as well. A common source for producing alcohols is from carbonyl compounds. The choice of carbonyl type (ketone, aldehyde, ester, etc.) and the type of reaction (Grignard addition or Reduction), will determine the product(s) you will get. Fortunately, there are a number of variations of carbonyls, leading to a number of choices in product. There are primarily two types of reactions used to create alcohols from carbonyls: Grignard Addition reactions and Reduction reactions. We'll look at each type of reaction for each type of carbonyl. Grignard Addition Reactions. As we learned previously, Grignard reagents are created by reacting magnesium metal with an alkyl halide (aka haloalkanes). The magnesium atom then gets between the alkyl group and the halogen atom with the general reaction: R-X + Mg → R-Mg-X In our examples, we'll be using bromine in our Grignard reagents because it's a common Grignard halogen and it will keep our examples a little clearer without the need for X. The general mechanism of a Grignard reagent reacting with a carbonyl (except esters) involves the creation of a 6-membered ring transition state. The pi bond of the oxygen attacks a neighboring magnesium bromide which in turn, releases from its R group leaving a carbocation. At the same time, the magnesium bromide ion from another Grignard molecule is attacked by the carbocation and has its magnesium bromide ion stolen (restoring it to its original state as a Grignard reagent). The second molecule's carbocation is then free to attack the carbanion resulting from the vacating pi bond, attaching the R group to the carbonyl. According to youtube, see Grignard reagent, the carbanion (R:-) from the R-MgBr attacks the partially positive carbonyl carbon, displacing the pi electrons onto the O, which takes a proton, forming the alcohol. At this point, there is a magnesium bromide on the oxygen of what was a carbonyl. The proton from the acidic solvent easily displaces this magnesium bromide ion and protonates the oxygen, creating a primary alcohol with formaldehyde, a secondary alcohol with an aldehyde and a tertiary alcohol with a ketone. With esters, the mechanism is slightly different. Two moles of Grignard are required for each mole of the ester. Initially, the pi bond on the carbonyl oxygen attacks the magnesium bromide ion. This opens up the carbon for attack from the R group of the Grignard. This part of the reaction is slow because of the dual oxygens off of the carbon providing some resonance stabilization. The oxygen's pi bond then re-forms, expelling the O-R group of the ester which then joins with the magnesium bromide, leaving R-O-MgBr and a ketone. The R-O-MgBr is quickly protonated from the acidic solution and the ketone is then attacked by Grignard reagent via the mechanism described earlier. Synthesis from Formaldehyde. The image above shows the synthesis of an alcohol from formaldehyde reacted with a Grignard reagent. When a formaldehyde is the target of the Grignard's attack, the result is a primary alcohol. Synthesis from an Aldehyde. The image above shows the synthesis of an alcohol from an aldehyde reacted with a Grignard reagent. When an aldehyde is the target of the Grignard's attack, the result is a secondary alcohol. Synthesis from a Ketone. The image above shows the synthesis of an alcohol from a ketone reacted with a Grignard reagent. When a ketone is the target of the Grignard's attack, the result is a tertiary alcohol. Synthesis from an Ester. The image above shows the synthesis of an alcohol from an ester reacted with a Grignard reagent. When an ester is the target of the Grignard's attack, the result is a tertiary alcohol and a primary alcohol. The primary alcohol is always from the -O-R portion of the ester and the tertiary alcohol is the other R groups of the ester combined with the R group from the Grignard reagent. Synthesis from an Epoxide. We will discuss reactions with Epoxides later when we cover epoxides, but for now, we'll briefly discuss the synthesis of an alcohol from an epoxide. The nature of the reaction is different than with the carbonyls, as might be expected. The reaction of Grignard reagents with epoxides is regioselective. The Grignard reagent attacks at the least substituted side of the carbon-oxygen bonds, if there is one. In this case, one carbon has 2 hydrogens and the other has 1, so the R group attacks the carbon with 2 hydrogens, breaking the bond with oxygen which is then protonated by the acidic solution. leaving a secondary alcohol and a concatenated carbon chain. The R group can be alkyl or aryl. Organolithium Alternative. As an alternative to Grignard reagents, organolithium reagents can be used as well. Organolithium reagents are slightly more reactive, but produce the same general results as Grignard reagents, including the synthesis from epoxides. IT IS KNOWN AS Organolithium Alternative. Reduction. Synthesis from an Aldehyde. The image above shows the synthesis of an alcohol from an aldehyde by reduction. Synthesis from a Ketone. The image above shows the synthesis of an alcohol from a ketone by reduction. Synthesis from an Ester. The image above shows the synthesis of an alcohol from an ester by reduction. Esters can be hydrolysed to form an alcohol and a carboxylic acid. Synthesis from a Carboxylic Acid. The image above shows the synthesis of an alcohol from a carboxylic acid reacted by reduction. =Properties= Naming alcohols. Follow these rules to name alcohols the IUPAC way: Acidity. In an O-H bond, the O steals the H's electron due to its electronegativity, and O can carry a negative charge (R-O-). This leads to deprotonation in which the nucleus of the H, a proton, leaves completely. This makes the -OH group (and alcohols) Bronsted acids. Alcohols are weak acids, even weaker than water. Ethanol has a pKa of 15.9 compared to water's pKa of 15.7. The larger the alcohol molecule, the weaker an acid it is. On the other hand, alcohols are also weakly basic. This may seem to be contradictory--how can a substance be both an acid and a base? However, substances exist that can be an acid or a base depending on the circumstances. Such a compound is said to be amphoteric or amphiprotic. As a Bronsted base, the oxygen atom in the -OH group can accept a proton (hydrogen ion.) This results in a positively-charged species known as an oxonium ion. Oxonium ions have the general formula ROH2+, where R is any alkyl group that is a carbon containing species that ranges from -CH3. Alkoxides. When O becomes deprotonated, the result is an alkoxide. Alkoxides are anions. The names of alkoxides are based on the original molecule. (Ethanol=ethoxide, butanol=butoxide, etc.) Alkoxides are good nucleophiles due to the negative charge on the oxygen atom. Producing an alkoxide. R-OH -&gt; H+ + R-O- In this equation, R-O- is the alkoxide produced and is the conjugate base of R-OH Alcohols can be converted into alkoxides by reaction with a strong base (must be stronger than OH-) or reaction with metallic sodium or potassium. Alkoxides themselves are basic. The larger an alkoxide molecule is, the more basic it is. =Reactions= Conversion of alcohols to haloalkanes. Recall that haloalkanes can be converted to alcohols through nucleophilic substitution. This reaction proceeds because X (a halogen) is a good leaving group and OH- is a good nucleophile. OH, however, is a poor leaving group. To make the reverse reaction proceed, OH must become a good leaving group. This is done by protonating the OH, turning it into H2O+, which is a good leaving group. H+ must be present to do this. Therefore, the compounds that can react with alcohols to form haloalkanes are HBr, HCl, and HI. Just like the reverse reaction, this process can occur through SN2 (backside attack) or SN1 (carbocation intermediate) mechanisms. As stated in the haloalkane chapter, the two mechanisms look similar but the mechanism affects the rate of reaction and the stereochemistry of the product. Oxidation of alcohols. Oxidation in organic chemistry always involves either the addition of oxygen atoms (or other highly electronegative elements like sulphur or nitrogen) or the removal of hydrogen atoms. Whenever a molecule is oxidized, another molecule must be reduced. Therefore, these reactions require a compound that can be reduced. These compounds are usually inorganic. They are referred to as "oxidizing reagents". With regards to alcohol, oxidizing reagents can be strong or weak. Weak reagents are able to oxidize a primary alcohol group into a aldehyde group and a secondary alcohol into a ketone. Thus, the R-OH (alcohol) functional group becomes R=O (carbonyl) after a hydrogen atom is removed. Strong reagents will further oxidize the aldehyde into a carboxylic acid (COOH). Tertiary alcohols cannot be oxidized. An example of a strong oxidizing reagent is chromic acid (H2CrO4). Some examples of a weak oxidizing reagents are pyridinium chlorochromate (PCC) (C5H6NCrO3Cl) and pyridinium dichromate (PDC). 

Free-radical halogenation. Free-radical halogenation is "a way to do a substitution other than by nucleophilic substitution." "or, simplified:" H's come off as H+ not H-, making nucleophilic substitution impossible. We must first remove a H+ to perform the reaction. Free radicals as a rule violate the octet rule and are highly reactive, reacting even with strong bonds. "then" The reaction continues happening in a chain mechanism called a free radical chain reaction. Free radical chain reaction. This has three steps: Running free radical chain reactions with Chlorine and Bromine yield their respective end products in different proportions, a quality that can be useful in the laboratory. 

Radical stability. Alkyl groups vary in stability. CH3° is quite reactive as it violates the octet rule. However, CH3+ is also a charged moiety and is therefore even less stable and harder to form. 

Alkenes. Alkenes are organic molecules that contain at least one carbon-carbon double bond, which is referred to as unsaturation. Saturation is when all of the carbons have all single bonds. If the double bonds are just next to each other, they are called cummulated dienes. If there is only one C-C bond between two double bonds, it is called a conjugated diene. If more than one single bond intervenes it is an isolated diene. The double bonds in alkenes are higher in energy than single bonds and are more reactive. Molecules will favor single bonds over double or triple ones when given a chance. Bond dissociation energy is the energy needed to break a bond. Ethene is the smallest, simplest alkene. It is planar and the angle between its bonds is approximately 120°. The C-H bonds are stronger in ethene than in ethane because the π bond draws electron density away from from the carbon, which draws electron density away from the hydrogens. This makes ethene unreactive and gives the carbons sp2 character, which is more s-like and lower in energy. IUPAC alkene nomenclature. Then follow the rest of the rules for naming halogenoalkanes Alkenols. Alkenols are compounds containing both a double bond and an OH group. To name, find the longest carbon chain that contains both the C=C and the OH-. The OH group gets the higher priority (lower number). Oxygen takes precedence over carbon.  (IUPAC name) The "2" refers to the double bond and can go at the beginning of the name or right before the "en". Preparation of alkenes, or how to make double bonds.. 1°, 2°, 3° all work for E2 reactions (remember that the methyl group doesnt have a second carbon for an elimination to function). E1 reactions. Elimination reactions are one way to produce alkenes. Learn more about them here. Elimination vs Substitution ? There are ways to predict if a reaction will follow an elimination or a substitution mechanism or pathway. Tertiary halides favor E1, SN1. Reactions of alkenes. Bromine. Electrophilic addition reaction between Bromine and Ethene gas: Complete Combustion. Like other hydrocarbons alkenes will combust in excess air or oxygen provided that there is sufficient activation energy for the combustion reaction. The following reaction is the complete combustion of ethene with oxygen: Hydrogen Bromide. Electrophilic addition reaction between concentrated Hydrogen Bromide and Ethene gas: The general formula for the reaction of any alkene with HBr is: For asymmetrical alkenes such as But-1-ene the following reaction is also feasible: Where R and R' are alkyl groups e.g. CH3 or CH3CH2 Sulphuric Acid. Electrophilic addition reaction between cold concentrated Sulphuric Acid and Ethene gas: Ethyl Hydrogensulphate is hydrolysed if warmed with water: Overall reaction for the two steps: « Alcohols | Alkenes | Alkynes » 

Electrophilic additions are also referred to as addition reactions. Electrophilic additions are essentially the reverse of an E1 elimination reaction, sometimes exactly the microscopic reverse. Addition reactions involve carbocations. This intermediate carbocation can rearrange. Addition reactions follow Markovnikov's rule. (The higher priority substituent adds to the more highly substituted carbon of the carbon-carbon double bond). 

Lessons. Punto de partida.  "A starting point" — the basics of conversation and grammar. ¿Qué opinas?  "What do you think?" — a world of opinions and viewpoints. ¿Cuándo se hará?  "When will it be done?" — The busy world of time. Other materials. The materials below are mostly miscellaneous or incomplete. Contributing to this Wikibook. This is a "Wiki"book—please feel free to edit, enhance, update and add to it, in any way that will make it a better teaching resource! Especially if you are a native Spanish speaker, please consider recording sound files for any of the dialogue or vocabulary. It seriously improves the learning potential of the book if you can listen as you read! A list of the people who have contributed to this featured book is here. The book already has all of the most basic information allowing you to construct a simple sentence, and say a few things in Spanish. However, do try to help by adding some more information on other topics that are not here. As you write, remember that people reading this book can be totally new to the language. Do not write as if you are writing for another Spanish speaker, but keep in mind to explain things in great detail to enable learners to use this book easily. Contribute to this book to make it a good way for new learners to learn the Spanish language! 

External links. /Grammar 

Formation. Most French verbs fall into the category of "-er" verbs. To conjugate, drop the "-er" to find the stem or root. Add endings to the root based on the subject and tense. Pronunciation, elision and liaison. The "-e", "-es", and "-ent" endings all have the same silent pronunciation. The "-er" and "-ez" endings are pronounced , and the "-ons" ending is pronounced . In all conjugations, "je" changes to "j" ' when followed by a vowel or silent "h": In all plural forms, the "s" at the end of each subject pronoun, normally unpronounced, becomes a "z" sound and the "n" of "on" becomes pronounced when followed by a vowel.  The direct object pronoun , meaning "it" or "him", replaces masculine singular direct objects: , meaning "it" or "her", replaces feminine singular direct objects: , meaning "them", replaces plural direct objects: "Le" and "la" become "l'" before a vowel: "Le", "la", and "les" can replace either people or things: The verb is a regular "-er" verb meaning "to play". It can be used to refer to sports, games, and instruments. When referring to sports or games, "jouer à …" is used; recall that replaces "à le", and replaces "à les": When referring to instruments, "jouer de …" is used; recall that replaces "de le", and replaces "de les":   The indirect object pronoun means "to him" or "to her": Likewise, the indirect object pronoun means "to them": "Lui" replaces "à ["person"]": Likewise, "leur" replaces "à ["people"]": "Lui" and "leur" usually only refer to people; they will sometimes also be used in reference to things. Examples. Like in English, some verbs can be followed by infinitives. Three "-er" verbs used in this manner are "aimer", "adorer", and "détester": Examples. The verb "s'amuser" means "to have fun" in English. It is a type of pronominal verb (a verb that includes a pronoun as part of it) called a reflexive verb, which means that the action of the verb is reflected back onto the subject. Literally translated, the verb means "To amuse oneself." Examples.  

Puzzles | Logic puzzles 

Puzzles | Logic puzzles | Knights, Knaves &amp; Spies On the fabled Island of Knights and Knaves, we meet three people, A, B, and C, one of whom is a knight, one a knave, and one a spy. The knight always tells the truth, the knave always lies, and the spy can either lie or tell the truth. A says: "C is a knave."&lt;br&gt; B says: "A is a knight."&lt;br&gt; C says: "I am the spy." "Who is the knight, who the knave, and who the spy?" Solution 

Puzzles | Logic puzzles | Knights, Knaves &amp; Spies II We have three people one of whom is a knight, one a knave, and one a spy. The knight always tells the truth, the knave always lies, and the spy can either lie or tell the truth. The three persons are brought before a judge who wants to identify the spy. A says: "I am not a spy."&lt;br&gt; B says: "I am a spy."&lt;br&gt; Now C is in fact the spy. The judge asks him: "Is B really a spy?" "Can C give an answer so that he doesn't convict himself as a spy?" Solution 

Puzzles | Logic puzzles | Lying about your Age Annie, Betty, Carrie, Darla, and Eve recently found out that all of their birthdays were on the same day, though they are different ages. On their mutual birthday, they were jabbering away, flapping their gums about their recent discovery. And, lucky me, I was there. Some of the things that I overheard were... Since I knew these people -- and how old they were, I knew that they were not telling the whole truth. After thinking about it, I realized that when one of them spoke to somebody older than herself, everything she said was true, but when speaking to somebody younger, everything she said was false. "How old is each person?" Solution 

Puzzles | Logic puzzles | Getting out of Prison You are in a prison and your execution is planned for today. However, you are given one last chance to get out of this situation. You are brought into a room that has two exits and in front of each of them stands a guard. One exit leads out of the prison into freedom, while the other one brings you into the room where the execution takes place. Furthermore, one of the guards is always lying, while the other one is always telling the truth (but, of course, you don't know who of the guards is in front of which door). You are allowed to ask only one question. "What question can you ask in order to find out the way to your freedom?" Solution 

Puzzles | Logic puzzles | Hats 3 men are buried in the sand all facing the same way with only their heads above ground. Each man has a hat placed on his head which is either red or blue, but they can't see it. The men know that their hats were selected from a bag containing 3 red hats, and 2 blue hats. The man at the back (who can see the hats of the two men in front of him) is asked what hat he is wearing. He answers: "I don't know". The middle (who can see the hat of the man in front of him) man is asked what hat he is wearing. He also replies: "I don't know". Finally, the man at the front is asked what hat he is wearing. He answers: "I am wearing a red hat". "How did he know?" Solution 

Les pays du monde (nations of the world) W. [None] X. [None] 



Puzzles | Geometric puzzles | Connecting Stars     "In the above picture cross all the stars with 4 straight lines and without taking the pencil off the paper (or screen;)!" solution 

Puzzles | Geometric puzzles | Rubber Band Imagine a rubber band which is infinitely thin and fixed on one end onto a plane with a nail. In its relaxed state it measures exactly 6cm and can be stretched to exactly 8cm. "Your task is to construct a right-angled triangle in the plane using only the rubber band, your absolutely precise hands and a sharp pencil." Solution 1 | Solution 2 



Puzzles | Arithmetical puzzles | Two 4's equals 64? You can use the number 4 twice (but no other numbers) and as many mathematical operators as you wish to write an expression that is equal to 64. "What is the expression?" solution 

Puzzles | Arithmetical puzzles | Three Daughters A salesman is at the door of a family house and tries to convince the mother to buy some of his products. She says: "Well, I don't really need the products, but if you can guess how old my three daughters are I will do you a favor and buy one of them." The salesman asks the woman for a little help and she tells him: "The product of their ages is 36." As the salesman is not able to figure out the ages, she gives him another clue: "The sum of their ages is same as the number of my house." The salesman calculates for a little while but then says: "I'm sorry, Miss, but I still can't figure out their ages." The mother agrees to give him a final hint: "My eldest daughter plays piano." After that hint the salesman was able to tell the mother her daughters ages and sell one of his items. "How did he figure that out?" solution 



Puzzles | Decision puzzles | 12 coins You are handed 12 metal coins which look alike. However, one of them is fake and weighs either more or less than each of the other 11 coins (which all have identical weight). To measure or compare weights you have only a pan balance that tips down to the side with the greater weight or stays in balance if both pans hold equal weights. "With how many weightings can you find the fake coin and also determine if its weight is more or less than the weight of the other coins? How do you do it?" solution 

Puzzles | Decision puzzles | Weighings once more You want to use a balance to weigh items that have an integral weight between 1 ounce and 40 ounces (integral means that no fractions such as 1.5 ounces are allowed). "What is the fewest number of weights you need?" solution 

Puzzles | Action sequences 

Puzzles | Action sequences | Sharing Milk You have an 8 liter bucket full of milk and two empty buckets, one with a 5 liter volume and the other with a 3 liter volume. You have no means of measurement other than the buckets themselves, which are not calibrated. "How can you evenly distribute the milk into two portions of precisely 4 liters?" Solution 

Puzzles | Action sequences | Crossing the River A farmer wants to cross a river and take with him a wolf, a goat, and a cabbage. There is a boat that can fit himself plus either the wolf, the goat, or the cabbage. If the wolf and the goat are alone on one shore, the wolf will eat the goat. If the goat and the cabbage are alone on the shore, the goat will eat the cabbage. "How can the farmer bring the wolf, the goat, and the cabbage across the river?" Solution 

Puzzles | Action sequences | Crossing the River II There are six persons on the shore: Three fathers and each of them has a son. The sons are afraid to be nearby a different father (or fathers) if their own father is not there (both at the boat and the either of shores). There is a boat that can hold a maximum of two persons (you can assume that the sons are able to row the boat). "How can the six people get across the river?" Solution 

Puzzles | How do you ... ? 

Puzzles | Puzzles/How do you ... ? | The Burning Island A man is stranded on an island covered in forest. One day, when the wind is blowing from the west, lightning strikes the west end of the island and sets fire to the forest. The fire is very violent, burns everything in its path, and without intervention the fire will burn the whole island, killing the man in the process. There are cliffs around the island, so he cannot get back on if he jumps into the ocean. "How can the man survive? (The fire, at least.) There are no buckets or any other means to douse the fire with water or sand to put it out." solution 

Puzzles | Arithmetical puzzles | Digits of the square There is one four-digit whole number "n", such that the last four digits of "n"² are in fact the original number "n". "What is it?" solution 

Puzzles | Decision puzzles | Monty Hall Let us assume that you are participating on the show "Let´s make a deal". Monty Hall offers you a choice of three doors. Monty tells you, and you know, that you can trust Monty, that there is one big prize (such as a sports car) behind one of the doors and goats behind the others. He tells you that after you choose a door, he will open a different door than you chose, which contains a goat. After you have chosen one door, he opens one of the two other doors, which as promised, contains a goat. Now he gives you the chance to switch or stay at your initial choice. You will then get what is behind that door. You cannot hear the goats from behind the doors, or in any way know which door has the prize. "Does it matter whether you switch or stay?" solution 

The task is to find the hidden principle that constructs the list, and sometimes to continue the sequence, which is indicated by a question mark. Be aware: some sequences are infinite, some are finite. If you find a different solution than the given one please add it to the solution page. Creativity is encouraged! Easy sequences.  1 2 4 8 ? solution  1 2 4 7 ? solutions (more than one)  3 9 27 81 ? solution  2 3 5 8 13 21 ? solution  1 0 1 0 1 0 1 1 0 1 0 1 solution  2 3 6 7 1 9 4 5 8 solution  0 1 1 1 1 1 1 ... solution Medium Sequences.  1 3 7 61 ? solution  2 -2 0 2 -3 0 ? solution  1 4 1 5 9 ? solution Hard Sequences.  85 90  29 67 0 solution  83 33 69 2  1 81 91 57 36  60 55 61 29 35 6  29 42 45 23 13 5 16  23 13 44 7 39 69 47 98  61 41 91 82 1 37 27 71 36  6 4 6 9 0 8 9 5 9 6  73 94  69 81 28 solution  56 3 43 54  33 98 49 91 8  7 31 44 27 57 6  26 45 71 38 61 65 88  25 36 7 97 76 91 87 86  83 55 98 1 5 10 1 15 96  3 9 3 4 4 3 0 7 5 6 

"-ger" verbs are regular "-er" verbs that are also stem changing. The most common "-ger" verb is "manger." For "manger" and all other regular "-ger" verbs, the stem change is the addition of an "e" after the "g". This only applies in the "nous" form. In this case, the change is made to preserve the soft "g" pronunciation rather than the hard "g" that would be present if the "e" were not included. The verb "boire", meaning "to drink", is irregularly conjugated (it is not a regular "-re" verb). The partitive article "de" indicates, among other things, the word "some". As learned earlier, "de" and "le" contract (combine) into "du", as "de" and "les" contract into "des". Also, instead of "du" or "de la", "de l"' is used in front of vowels. When speaking about food, the partitive article is used at some times while the definite article ("le", "la", "les") is used at other times, and the indefinite article ("un", "une") in yet another set of situations. In general "de" refers to a "part" of food (a "piece" of pie) whereas the definite article ("le") refers to a food in general ("I like pie (in general)"). The indefinite article refers to an "entire unit" of a food ("I would like a (whole) pie"). When speaking about preferences, use the definite article: When speaking about eating or drinking an item, there are specific situations for the use of each article. In the negative construction, certain rules apply. "Un" or "une" changes to "de" in a negative construction, meaning, in this context, "any". Similarly, "du", "de la", or "des" change to "de" in negative constructions. To say "some of it" without specifying the exact object, the pronoun "en" can be used. Additionally, "en" can mean "of it" when "it" is not specified. For instance, instead of saying "J'ai besoin d'argent", if the idea of money has already been raised, it can be stated as "J'en ai besoin." This is because "en" replaces "du", "de la" or "des" when the noun is not specifically mentioned in the sentence. Like with "me", "te" and other pronouns, "en" (meaning "some") comes before the verb. 

"Jean and Chantal are discussing what types of pen they have." Jean - J'ai un stylo rouge. Chantal - Moi, j'ai un stylo bleu. Et toi, tu as un stylo rouge. Jean - Aussi, j'ai un stylo jaune. "Of course, you can use avoir for anything you have!" Chantal - J'ai deux frères et trois sœurs, et toi ? Jean - J'ai un frère et une tante. Chantal - J'ai deux tantes et un oncle. "Avoir is also used to describe age" Jean - J'ai quatre ans. Et toi ? Chantal - J'ai trois ans. "Avoir", meaning "to have", is conjugated irregularly. Formation. Remember to "liaison" between "nous avons", "vous avez", and "ils ont/elles ont". Expressing age. "Avoir" is used to express age. Interrogatives. The above uses avoir "affirmatively". You can also use it "interrogatively". A small complication arises, in that without some help, the result does not sound very good. The use of an "euphonic" (pleasing to the ear) is used with vowels before the pronoun. Thus, the letter "-t-" is placed between the verb and the pronoun: Ai-je ? (Have I ?) As-tu ? (Have you ?) informal (hast thou) A-t-il ? (Has he ?) A-t-elle ? (Has she ?) Avons nous ? (Have we ?) Avez vous ? (Have you ?) formal Ont ils ? (Have they ?) masculine Ont elles ? (Have they ?) feminine The use of "liaison" fullfils the euphonic for "ont". Examples: A-t-il la farine ? Oui, Monsieur, il a la farine. Avons nous la viande ? Oui, Monsieur, nous avons la viande et le pain. Avez vous la table ? Oui, Madame, j'ai la table. The preposition is used to express possession or association: "De" can also be translated as "'s": Recall that replaces "de le", and replaces "de les". Usage. Possessive adjective are used to express possession of an object. In English the possessive adjective agrees with the subject ("his sister", "her brother"). But in French, possessive adjectives act like all other adjectives: they must agree with the noun they modify. Whether "son", "sa" and "ses" translate to "his" or "her" is indicated by context: "Notre", "votre", and "leur" modify singular nouns, regardless of gender; "nos", "vos", and "leurs" modify plural nouns: Liaison and adjective changes. Liaison occurs when "mon", "ton", and "son" are followed by a vowel. Liaison also occurs with all plural forms, since they all end in "s". "Mon", "ton", and "son" are used before a feminine singular noun that starts with a vowel or silent "h": Examples.  Recall that the expression means "there is …" or "there are …": The interrogative form of "il y a" is . That is, "il y a" is inverted to "y a-t-il", meaning "is there?" or "are there?", within questions: Both "How much …" and "How many …" are translated as . If the person or thing it refers to is countable, "combien de" is always followed by a plural noun: However, with uncountable nouns, such as and , the singular form is used: As with "il y a", other nouns and verbs can be inverted within questions. For example, can become : Examples. To speak about complex family relations, you use "de mon", "de ma", and "de mes": Examples.  

Direct objects. While the subject of a sentence initiates an action (the verb), the direct object is the one that is affected by the action. A direct object pronoun is used to refer to the direct object of a previous sentence: The following table shows the various types of direct object pronouns: Notes: Indirect objects. An indirect object is an object that would be asked for with "To whom...?" or "From whom...?". It is called indirect because it occurs usually together with a direct object which is affected directly by the action: The following table shows the various types of indirect object pronouns: Notes: The bread "is given" by the man (direct). Pierre "gets the given" bread (indirect). Example : "-e…er" are regular "-er" verbs, but also are stem changing. The stem change applies to all forms except "nous" and "vous". The stem change involves adding a grave accent ( ` ) over the "e" in the stem. Formation.  "-yer" verbs are irregular "-er" verbs. When "y" is part of the last syllable, it changes to "i" in order to keep the "ay" sound. In the present indicative of "-yer" verbs, this affects all forms except "nous" and "vous". Some "-yer" verbs, such as "payer", may optionally retain the "y". Formation. In the present indicative, is "conjugated" as follows:  Many of the verbs you have learned so far have irregular past participles.   

 The is a compound tense, and is therefore composed of an auxiliary verb and a past participle. With most verbs, that auxiliary verb is "avoir". Meaning. In English, verbs conjugated in the "passé composé" literally mean "have/has ____ed". While there is a simple past tense in French, it is almost always only used in formal writing, so verbs conjugated in the "passé composé" can also be used to mean the English simple tense. For example, the "passé composé" forms of , "["avoir"] parlé", literally mean "has/have spoken", but also means "spoke". Basic formation. To conjugate a verb in the "passé composé", the helping verb, usually "avoir", is conjugated in the present indicative and the past participle is then added. Auxiliary verb - "avoir". Conjugate "avoir" in the present indicative. "Avoir" + past participle. Please also note: Fem. Subject or Person (Elles, Elle, Nous, On etc.)- Add another e with no aigu or grave to end of word- if a female person is partaking in the group. Plural Subject (On, Nous, Tu, Vous etc.)- Add another "s" to end of word. Finally, some verbs are irregular for the past participle, such as aller (to go), instead of using avoir to form the past participle, they will use être (to be) to form the past participle. Always check the verb's irregularities before using to form past participle. Some "past participle" irregulars are regular verbs when forming other tenses. Examples.  The word "professeur" is considered masculine at all times, even if the teacher is female. The only case when "professeur" can be preceded by feminine determinant is either when contracting it in colloquial language "la prof", or when adding a few words before : "madame/mademoiselle la/le professeur". In French, you do not "own" body parts. While in English, you would say "my hand" or "your hand", the definite article is almost always used in French: "To" and "of" are part of the verbs "écouter" and "entendre" respectively. It is not necessary to add a preposition to the verb. Other verbs, such as "répondre (à)", meaning "to respond (to)", are almost always followed by a preposition. Écrire. "Écrire" is an irregular French verb, meaning "to write". It varies from other "-re" verbs in the plural conjugation, by adding a "v". Its past particple, "écrit", is also irregular. The verb is conjugated the same way. The nouns , meaning "writing" or "handwriting", and , meaning "writer", are derived from "écrire". Lire. "Lire" is an irregular French verb, meaning to read. Its plural conjugation adds an "s", and its past participle is "lu". The verbs and are conjugated the same way. The adjective , meaning "readable" or "legible", is derived from "lire". The way that grades are numbered in France is opposite the way they are in the US. Whereas American grade numbers increase as you approach your senior year, they descend in France. 

Usage of the Subjunctive Mood. The subjunctive mood is used to express subjectivity, as opposed to objectivity. The subjunctive usually appears with "que", which means "that". 



Note: The indicative indicates certainty about an action. The subjunctive indicates a doubt or subjectivity. The conditional indicates that an action will occur or occurred based on the fulfillment of certain conditions. 

Usage of the Future. One uses the future tense when referring to an action, certain to occur, in the future. In a time ahead of now. One may also use "aller" in the present tense in conjunction with "aller" or another verb in infinitive form, to refer to the future. However it is not the future tense. For example, "Il va aller à l'école" Or "Je vais dormir" Holds generally the same meaning as, "Il ira à l'école" Or "Je dormirai" However, the former is not in the future tense. Also, the usage of "aller" generally signifies an action to occur in the very near future, where as future tense refers to any time in the future. Formation of the Future. Future Stems. Stem Changes (Infinitif) Exceptions to the Rule (Irrégulier). Future Endings. To conjugate a verb in the futur simple, one takes the infinitive and appends the following, as according to the table: 

Usage of the conditional (present). The conditional tense is used when: Formation of the conditional. Conditional stems. Stem changes Irregular stems 

"Imparfait" in French Usage of the Imperfect. The imperfect is used in French under several different circumstances. The imperfect is used: Formation of the Imperfect. Examples. Aller Être Exceptions. Manger Commencer Exercises. Translate the following sentences into English: -Je "jouais" au foot quand j'"avais" douze ans, mais maintenant je nage parfois. Quand j'"avais" douze ans j'"étais" en forme. Une fois, le douze décembre, je "me suis cassé" la jambe, et je ne "jouais" plus au foot. Quelle tristesse! -Quand j'"avais" dix ans, je "mangeais" beaucoup de frites. 

Usage of the Simple Past (Past Historic, Past Definite). The simple past is mostly a literary tense, used in fairy tales, and perhaps newspapers. It is one that native French students are expected to recognize but not use. Formation of the Simple Past. To conjugate in this tense, one finds the stem and appends the following, as according to the table: It should be noted that être, along with a few other verbs are consistent in their irregularities in the passé simple as well. Simple Past Stems. -er, changes it to é (manger = mangé) &lt;br&gt;-ir, take off the r (choisir = choisi) &lt;br&gt;-re, take off the re and add a u 

The passé composé is a perfect tense, and is therefore composed of an auxiliary verb and a past participle. With most verbs, that auxililary verb is avoir. Meaning. Verbs in French conjugated in the passé composé can most simply be translated to English as eg "has / have ____ed". While there is a simple past tense in French, it is mostly used in formal narrative writing, so verbs conjugated in the passé composé can also be used to mean the English simple tense. When to use. You use the passé composé when you want to express that: Formation. Introduction. To conjugate a verb in the passé composé, the auxiliary (or helping) verb, usually avoir, is conjugated in the present indicative and the past participle is then added. It is important to remember that there is only *one* verb in the passé composé. While the past participle looks like a verb, it is not - it functions more like an adjective. This is important to remember because when you negate in the passé composé, you negate the "only" verb, which is the auxiliary verb (ex. "Je n'ai pas mangé"; "I have not eaten"). This works exactly the same way in English - the only verb is the auxiliary verb, which is also the only thing negated in English ("I "have not" eaten"). Formation Summary. The compound past is a compound tense- it consists of two verbs, the auxiliary verb ("helper verb") and the past participle of the verb one seeks to use in this tense. To form the passé composé, you need to take the auxiliary verb - either avoir or être, then conjugate it according to the subject of the sentence, like in the present indicative tense. We then take the past participle of the verb, and stick that on the end. Every verb has one past participle that does not change (there are some exceptions, as one will learn later). To find the past participle, the stem of the infinitive must be determined or the irregularity must be known. If we want to make the statement negative, for example if we didn't do something in the past, we must always put the negative structure such as "ne ... pas" around the auxiliary verb, immediately before the past participle. For example, ""Je ne peux pas",". Also, reflexive or pronomial verbs must be conjugated with être under most circumstances. For example, the verb "se réflechir" is conjugated in the first person singular by "Je me suis réflechi(e),". Auxiliary Verb Formation. Auxiliary Verb - Avoir. Conjugate avoir in the present indicative. Auxiliary Verb - Être. Conjugate être in the present indicative. Past Participle Agreement with Preceding Direct Objects. The past participle must agree with the direct object of a clause in gender and plurality if the direct object goes before the verb. Avoir ou Être? In most circumstances, the auxiliary verb is avoir. However, with certain verbs, the auxiliary verb is être. This occurs under two different circumstances: 1. Reflexive verbs always take être. 2. The House of Être: Most verbs form the "passé composé" with "avoir", however there are a small number of verbs that are always conjugated with "être". Seventeen special intransitive verbs take "être" (four of which can also take "avoir", as explained below). 2.a. Exceptions Note that there are four verbs above that are followed by a star ("sortir, descendre, monter, passer"). When a direct object is used with these verbs, the auxiliary verb becomes "avoir". 

 Lektion 1 Grammatik 1-1 ~ Introduction to German grammar. Knowing the parts of speech (how words function in a sentence) is important for anyone attempting to learn a second language. English speakers will find many strong parallels between their language and German. However, as noted in the introduction, German grammar signals—how words indicate their function in a sentence—are more complex than English, and identifying the meaning of words in a German sentence is difficult without understanding these clues or signals to word function that come from the grammatical rules. The basic lessons (Level II) of this textbook are set up to first introduce the parts of speech, and then bring in the rules that govern these. Pay particular attention to both word endings and sentence word order as you progress in learning the German language. Following is a short conversation piece ("Gespräch"). Play the audio file first, then attempt to repeat what you hear, reading the spoken parts of the conversation. Go back and forth (listening and then speaking) until the German flows easily from your lips. This may take considerable practice. Refer to the vocabulary ("Vokabeln") below to understand the meaning of the German sentences you are hearing and speaking. Gespräch 1-1 ~ "Die Freunde". In this conversation we learn several simple greetings exchanged between friends meeting very briefly on the street. Vokabeln 1-1. This first vocabulary ("Vokabeln") may seem a bit long considering you have been presented with only the brief conversation piece above, but it also contains all of the German words you have encountered up to this point in the Level II textbook, including words in photo captions and lesson section headers. The layout of the "Vokabeln" is explained in the Lesson Layout Guide in the German~English textbook introduction, but the four parts of the "Vokabeln" are labeled in this first lesson to reenforce the concept. Note that column 3 may contain (in parentheses) additional notes about a word in column 1. Also, you can find the greeting phrases that appear in the simple conversations above (and many others) in Appendix 2, a German-English phrase book.  NOUNS  der Anhang, die Anhänge appendix, appendices (singular and plural)  die Brücke bridge  der Freund, die Freunde friend, friends (singular and plural)  das Gespräch, die Gespräche conversation, conversations  die Grammatik grammar (note irregular stress)  die Lektion lesson (note irregular stress)  die Straße street  das Tor gateway  die Vokabeln word list, vocabulary  das Vorwort foreword, preface (introduction to a book)  SHORT PHRASES  auf der Straße on the street  Auf Wiedersehen Good bye  Mir geht es gut I am fine (lit: 'It goes with me good')  Guten Tag! Good day (greeting)  Und dir? And you? (implied: 'And how are you?')  unter Freunden between friends  Wie geht es dir? How are you (lit: 'How goes it with you?')  Wie geht's? How are you? (casual, but more commonly used)  VERBS  gehen go ("geht" is "goes")  treffen meet, come upon ("trifft" is "meets")  OTHER "SMALL" WORDS (adjectives, adverbs, prepositions, etc.)  danke thank you; thanks  dir (with or for) you  einfach simple  es it  gut good  mir (with or to) me  und and  wie? how? Gespräch 1-2 ~ "Die Studenten". Here again, two friends (college students) meet casually and discuss briefly what each is doing. Grammatik 1-2 ~ Word Order in Questions. Basic or normal word order in simple German sentences is the same as in English—subject then verb then verb object: Unlike with English sentence structure, a question sentence in German is formed by reversing subject and verb: This is called inverted word order. Examples are provided in Gespräch 1-1 and Gespräch 1-2. As another example, consider the statement: "Er studiert Biologie" ('He studies biology'). A question statement might be: "Was studiert er?" ('What studies he?'; although in English, we would usually say: "What is he studying?"). The normal word order of subject ("er" or "he") then verb ("studiert" or "study") is reversed and, in this case, an interrogative ("was" or "what") added onto the front replacing the unknown (to the speaker) object (here, "biology"). Additional examples of questions formed from basic statements illustrate inverted word order: Grammatik 1-3 ~ Introduction to pronouns. A ("Pronomen") is a short word that takes the place of a noun previously mentioned in the sentence, paragraph, or conversation. A pronoun substitutes for a noun or noun phrase and designates persons or things asked for, previously specified, or understood from context. A specific pronoun in English as well as German has person, number, and case. You will be encountering all of the common German pronouns in the next several lessons, so we will track these as they appear. The following familiar personal pronouns are introduced in this lesson ("Lektion 1"):  "ich" – I (1st person, singular, nominative case)  "mich" – me (1st person, singular, accusative case)  "mir" – me (1st person singular, dative case)  "du" – you (2nd person, singular, nominative case)  "dich" – you (2nd person, singular, accusative case)  "dir" – you (2nd person singular, dative case)  "er" – he (3rd person singular, nominative case)  "sie" – she (3rd person singular, nominative case)  "es" – it (3rd person singular, nominative case) Pronoun person describes the relationship of the word to the speaker (that is, "1st person" is the speaker; "2nd person" is spoken to; and "3rd person" is spoken about). Pronoun number refers to whether the word represents one ("singular") or more than one ("plural") person or object. Finally, case indicates how the pronoun is used in a sentence, as will be explained over the next several lessons. For now, note in the examples you have already encountered, the three cases of 1st person singular pronouns in German: "ich", "mich", and "mir". In English these are: 'I', 'me', and ("to" or "with") 'me' — in essence, there are really just two cases in English: subjective ('I') and objective ('me'). You will shortly see that there are similarities, yet distinct differences, in the cases as used by the English and German languages. Vokabeln 1-2.  NOUNS  die Antwort, die Antworten answer(s) (singular and plural)  die Biologie biology (note irregular stress)  die Freundin, die Freunde (female) friend, friends (compare "der Freund")  der Käse cheese  der Kühlschrank refrigerator  die Mathematik mathematics (note irregular stress)  das Pronomen pronoun (note irregular stress)  der Student, die Studentin student, (female) student  die Uni university (a short form of "die Universität")  die Übersetzung translation (lit. "over-setting")  die Universität university (note irregular stress)  die Wurst sausage, banger  SHORT PHRASES  Dann bis bald! then until (we) soon (meet again) ("until then")  zu tun to do  VERBS  begegnen meet  brauchen need, want, require  einkaufen gehen go shopping  haben have  studieren study  verstehen understand  OTHER "SMALL" WORDS  an to (towards)  bald soon  bis until  dann then  du you  er he  fast almost  hallo hello  ich I  leer empty, vacant  mich me  schön beautiful (in this case, 'nice' or 'fine')  sehr very  sie she  tschüss so long (good bye)  viel much  was? what?  wohin? where? Übersetzung 1-1. By referring back to lesson examples, you should be able to write out the following sentences in German. On a piece of paper, first number and write each English sentence. Then review the lesson above and produce a German sentence that says the same thing as each English sentence. After all seven lines are translated, follow the "Antworten" (answers) link to compare your work with the correct ones. Do not be too concerned at this point if your spelling of the German verbs do not match the answers. You will learn all about German verb forms in later lessons. 

Pathetic drawing following  / | \  ---- | ----  ------- Weight to be measured is #. We prove, that 4 is the minimum size of a weighing set to measure the numbers from 1 to 40.  Wl = Wr + x where Wl and Wr are the sum of all weights on left resp. right side. x has a unique solution for every pair of Wl and Wr. Now, we construct a minimum solution: Lemma: The weighing set S(n) = {1, 3, ..., 3n} can measure all weights between 1 and |S(n)|. Proof by induction: For S(0) the lemma trivially holds. Assume that the lemma holds for S(n). We will now prove that it holds for S(n+1). Because S(n) is a subset of S(n+1), we can surely measure all weighs between 1 and |S(n)| = (3n + 1)/2 without using the weight 3n+1. By adding this weight on the left side, we can reach all solutions between 3n+1 - (3n+1)/2 = |S(n)| + 1 and 3n+1 + (3n+1)/2 = 3n+2/2 = |S(n+1)|. QED. Because of |S(3)| = 40, a solution to the problem is the weighing set {1,3,9,27}, which is minimal. Questions: Algorithm for finding how to weigh a particular integer weight from 1 to S(n): 1 - Let x be the weight you want to find weighing set for it. 2 - Subtract x from S(n) and convert the result to base 3. Then zero pad to make the size of resulting string equal to the size of weighting set. 3 - Subtract one from each bit. Since the string is in base three, the result after subtraction is -1, 0, or 1. Reverse the sign of each number. 4 - Each digit determine where to use its corresponding positional weight in reverse order (i.e. 81, 27, 9, 3, 1). Zero means the corresponding weight not on the scale, +1 means on the right, and -1 on the left resp. As an example suppose we want to know how to place each of 5 weights (1, 3, 9, 27, 81) to weigh 51. S(n) = 121 121 - 51 = 70 base3(70) = 2121, zero pad =&gt; 02121, subtract one from each digit =&gt; -1, 1, 0, 1, 0, reverse the signs =&gt; 1, -1, 0, -1, 0 =&gt; (+1 x 81) + (-1 x 27) + (0 x 9) + (-1 x 3) + (0 x 1) = 51 These digits determines the placing of weights: 81 on right resp, 27 on left, 9 out, 3 on left, and 1 out 

Usage of the Past Conditional. Past conditional is used to refer to an event that could have taken place in the past. Eg. "If he had not got hungry, we would have gone further." 



Usage. This is added after an auxiliary verb for all the composed tenses: Formation. The table below shows additions to the normal past participle that must be made based on the gender and number of the subject. 





Usage. This is used in a sentence when there is something in a future tense, but this action is also in the future, but before the other future. This is called the "futur anterieur" in French. 

In French the "pluperfect" is called "le plus-que-parfait". In English, it is also called the "more than perfect". Usage. The pluperfect is used to describe a past action that occurred before a second past action that is in the "passé composé" or "imparfait". 

Usage of the Past Imperative. The past imperative is only ever used for giving commands one would like to have done - this is a rare literary mood as the present imperative is used more frequently. Chances are that you'll never need to know this mood in your life, let alone use direct, indirect pronouns and negations with this! 1) Aie écrit ce rapport demain - Have this report written tomorrow. &lt;br&gt; 2) Soyez partis à midi - Leave / be gone by noon.&lt;br&gt; 3) Ayons fini les devoirs à 7h00 - Let's have our homework finished by 7h00.&lt;br&gt; 4) N'aie jamais mangé tes légumes - Never have your vegetables eaten.&lt;br&gt; 5) Ne l'aie pas déclarée coupable - Have her not found guilty.&lt;br&gt; 

Usage of the Pluperfect Subjunctive. The French pluperfect subjunctive is the least common literary tense - it's the literary equivalent of the past subjunctive. Like all literary tenses, the pluperfect subjunctive is used only in literature, historical writings, and other very formal writing, so it is important to be able to recognize it but chances are that you will never in your life need to conjugate it. 



Usage of the Imperfect Subjunctive. The subjunctive imperfect is very rarely employed in French; generally it only appears in literature and is viewed as archaic. It can in all instances be replaced by the subjunctive present. The subjunctive imperfect is employed in any instance in which the subjunctive is required, provided the trigger verb is in a past tense. In the example "Il fallait que le garçon allât à l'école", the subjunctive trigger verb "falloir" is in the imperfect, thus "aller" has been conjugated in the subjunctive imperfect. French speakers would normally express this as "Il fallait que le garçon aille à l'école", where "aller" has been conjugated in the present subjunctive. Formation of the Imperfect Subjunctive. Imperfect Subjunctive Endings. Imperfect Subjunctive Endings Irregular Conjugations. Imperfect Subjunctive of Venir 

Usage of the Imperative. This tense is used to give commands, express requests or make suggestions. Formation of the Imperative. The imperative is used in "tu", "nous" and "vous" forms; the "nous" and "vous" forms are the same as the indicative in both regular and irregular verbs (except the 3 irregulars shown below). The "tu" form is also the same unless it comes from an infinitive that ends in -er, in which case the "tu" form would drop the 's' (e.g.: parles -&gt; parle). You could also drop the 's' when an -ir verb has the same endings as an er verb. The infinitive can also be used as the imperative, but only for impersonal commands, e.g.: "mettre la ceinture". Regular Conjugations. Danser Perdre finir Irregular Verbs in their Imperative Conjugations. Faire Aller Venir Sortir Irregular Conjugations. Être Avoir Savoir Vouloir 

Usage of the Present. When you want to talk about something that's happening now. Formation of the Present. To form the present tense, there are seven categories of verbs that you need to know about, sorted by their endings, and if they are regular (follow the rules) or irregular (have their own rules). They are: and Irregular Conjugations. Stem Changing -er Verbs. -cer Verbs -ger Verbs -ayer, -oyer, and -uyer Verbs -eler Verbs -eter Verbs Verbs with "e" in the second to last syllable Verbs with "é" in the second to last syllable Irregular -ir Verb Patterns. Couvrir Dormir Irregular -re Verb Patterns. -aindre, -eindre, and -oindre Verbs Completely Irregular Verbs. Aller - To Go S'asseoir - To Sit Avoir - To Have Boire - To Drink Conduire - To Drive Connaître - To Know (A person) Croire - To Believe Devoir - To Have To Dire - To Say/To Tell Écrire - TO Write Être - To Be Faillir - To Almost Do Faire - To Do/To Make Falloir - To Have To/To Need (This verb can only be used in the impersonnal form) Lire - To Read Mettre - To Put/To Place Mourir - To Die Naître - To Be Born Prendre - To Take Rire - To Laugh Savoir - To Know (a fact) Sortir - To Go Out Venir - To Come Vivre - To Live Voir - To See Vouloir - To Want 

 Lektion 3 "Die Zahlen" ~ The Numbers Lektion 3 ~ "Zählen von 1 bis 12". Counting in any language is a valuable skill best learned early on. In German as in English, there are both cardinal (counting) and ordinal (place or order) numbers, and number formation is similar in that the first twelve numbers are unique. Above twelve, numbers are formed by combination. For example, 13 is "dreizehn" and 14 is "vierzehn". Higher numbers will be the subject of later lessons. Note in the table how ordinals are formed from the cardinals in German by adding "te". 'Ten' becomes 'tenth' in English; "zehn" become "zehnte" in German. As in English, there are several nonconforming variants: "erste", "dritte", and "siebte". Audio: (385KB) Grammatik 3-1 ~ Telling time (hours). Knowing the numbers from 1 to 12, you can now begin asking and telling time in German. &lt;br&gt; Gespräch 3-1. Asking for the time is accomplished by the sentence: "Wie spät ist es?" ("How late is it?"). The answer places the hour in the line "Es ist ____ Uhr" ("It is __ o'clock"), substituting the correct cardinal value (except "ein" is used instead of "eins"). One could also ask: "Wieviel Uhr ist es?" (not used very often anymore) or respond "Es ist eins" or "Es ist drei", etc.—which may be imprecise, unless the time is close to the hour. The following sentences also relate to telling time: Knowing how to express the quarter, half, and three quarter hours will allow you to give the time more precisely. We will, of course, revisit this subject. Once you know how to count beyond twelve, the hour's division into 60 minutes can be expressed. Also, Germans (like most Europeans) utilize what is known in America as "military time" or a 24-hour clock. Vokabeln 3-1. Also included in the vocabulary for Lesson 3 are the ordinal and cardinal numbers 1 through 12 from "Lektion 3" above.  der Ball ball  der Junge, die Jungen boy, boys  das Lernen learning, study  der Nachmittag afternoon  die Stunde hour  die Uhr watch (timepiece); also "o'clock"  der Uhrturm clock tower  die Uhrzeit time, time of day  das Viertel quarter  die Zahl, die Zahlen number, numbers  bis zwei Uhr until two o'clock  das ist gut very well (lit.: "that is good")  eines Nachmittags one (unspecified) afternoon  ich kann... spielen I can play  es ist it is  willst du ...? do you want ...? (familiar form)  fragen ask (a question)  spielen play  zählen count  dann then  halb half, halfway to  nach about, after  spät late  vor before, until  zu to Grammatik 3-2 ~ Introduction to Nouns. A is a fundamental part of speech, occurring in sentences in two different ways: as subjects (performers of action), or objects (recipients of action). As a generality, a noun is the name of a "person, place, or thing". Nouns are classified into proper nouns (e.g. "Janet"), common nouns (e.g. "girl"), and pronouns (e.g. "she" and "which"). A proper noun (also called "proper name") is a noun which denotes a unique entity. The meaning of a proper noun, outside of what it references, is frequently arbitrary or irrelevant (for example, someone might be named Tiger Smith despite being neither a tiger nor a smith). Because of this, they are often not translated between languages, although they may be transliterated — for example, the German surname "Knödel" becomes "Knoedel" in English, as opposed to "Dumpling". Proper nouns are capitalized in English and all other languages that use the Latin alphabet; this is one way to recognize them. However, in German both proper and common nouns are capitalized (as are certain formal pronouns; see Grammatik 2-3). Grammatik 3-3 ~ Gender of Nouns. We have seen evidence of word gender in the pronouns we have been encountering; notably 'he', 'she', and 'it' in English and "er", "sie", and "es" in German. Just like many other languages (but not English), German has genders for nouns as well. Noun gender is indicated by the "definite article", which should always be learned as part of the noun. For this reason, nouns presented in each lesson's "Vokabeln" include the gender appropriate definite article. Definite Articles. The definite article ("bestimmter Artikel") is equivalent to an English 'the', and the three basic gender forms of definite articles in German are as follows: To say 'the book' in German, you would say "das Buch", because "Buch" is a neuter noun. To say 'the man' in German, you would say "der Mann", because "Mann" is a masculine noun. To say 'the woman' in German, you would say "die Frau", because "Frau" is a feminine noun. Noun gender does not always derive from actual gender where gender might be applicable. For example, 'the boy' is "der Junge" ("masculine"); but 'the girl' is "das Mädchen" ("neuter"). Also, nouns that have no inherent gender are not necessarily neuter. From this lesson: 'the watch or time piece' is "die Uhr" ('feminine'). Because German is generally more structured than English, it is important when learning German nouns to always learn them with their gender correct definite article; and in the "Vokabeln" nouns are always given with their associated definite article. That is, you must memorize the word for 'book' in German as "das Buch", not simply "Buch". Not just definite articles, but indefinite articles and adjectives have endings that must match the gender of the noun they precede. Using the wrong gender can alter the meaning of a German sentence, so in forming a proper sentence with "Buch", you will need to know that it is a neuter noun. Indefinite Articles. In addition to the definite articles—"the" in English and "der"-words in German—discussed above, both languages have indefinite articles ("unbestimmter Artikel"). Indefinite articles precede nouns in the same way that definite articles do, but convey a general or indefinite sense. These are "a" or "an" in English. Thus, 'the book' or "das Buch" refers to a definite or specific book, whereas 'a book' or "ein Buch" is indefinite about which book is referred to. Indefinite articles also have gender as shown here: Here are some examples of indefinite articles (underlined) used in German sentences: Why, you ask, are there words like "einen" in some sentences above—a spelling that does not appear in the gender table? The tables for both the definite and indefinite articles above are simplified at this stage, giving only articles in the nominative case (applied to words that are subjects of verbs). In the very next lesson you will start to address all the other cases in German. However, the nominative case is the one used to signify the gender of a noun, as in our "Vokabeln". Vokabeln 3-2.  das Buch book  die Frau woman  der Knödel dumpling  das Mädchen (young) girl  der Mann man  lesen read Übersetzung 3-1. Translate the following sentences into German: 



Presidents of the United States. Although Washington was a member of the Whig Party before the Revolution, after the war he was not a member of any party, though he tended to lean toward Federalist positions. Since the formation of the Democratic-Republican party and the Federalist Party, there has always been at least one viable political party. Today the United States has a two party system. There have been many third party movements, such as Ralph Nader, and Theodore Roosevelt, but these attempts to create a three-party system have, thus far, failed. 

This Wikibook begins with an introduction to the language and a series of lessons. You may also wish to browse further down this contents page for other useful pages, including Kosa kata (vocabulary), Alat-alat pembelajaran (Tools for learning) and other resources. Pelajaran "(Lessons)". Dasar "(Basic)". This section describes the basics of the Indonesian language. Mostly it discusses the most basic structure of the language. What you'll learn here will be clarified further in the later sections. Pemula "(Beginner)". This section will expand your vocabulary and grammar building on the sentence structures explained in the basic lessons. Ahli "(Expert)". Learn some cultural background of the Indonesian language and further improve your vocabulary using the grammar you've learned so far. Latihan "(Exercises)". Under construction Pranala Luar "(External links)". Kamus "(Dictionary)". Online dictionary links: 

Downloadable and Print Versions. If a lesson, grammar page, appendix, ot text has been added or the name of an existing page has been changed, please update the print version. History. Roby Joehanes initiated this wikibook in 2003. However, unhappy with the collaboration in Wikibook, he resigned in late 2003. This wikibook is then developed by Authors. The Indonesian textbook was started by Edit: If you look at Robby Jo's profile, it states that he has given up contributing to WikiBooks Indonesia. Although he will be sorely missed [from this place], his presence is much needed. Either we petition for him to come back, or we continue this great Wiki from where he left off ~Cheetincheetah Other Contributors: 

^ Indonesian ^ | « Lesson 0: The Alphabet | Lesson 1: Greetings |Lesson 2: This, That » Contoh Percakapan "(Dialogue Example)". Budi: Selamat pagi, Bu! Wati: Selamat pagi, Pak! Budi: Apa kabar? Wati: Baik. Anda? Budi: Baik-baik juga. Wati: Kamu sedang apa? Budi: Aku sedang membaca novel. Wati: Novel apa yang kamu baca? Budi: Aku sedang membaca novel semua tentang Islam! Wati: Boleh aku meminjamnya? Budi: Tentu saja. Wati: Sudah waktunya aku harus pulang. Sampai berjumpa, dan jangan lupa besok aku akan meminjam novelmu Budi: Tentu aku tidak akan lupa. Sampai jumpa kembali. Wati: Selamat jalan. Terjemahannya (The Translation): Budi: Good morning, Ma'am! Wati: Good morning, Sir! Budi: How are you? Wati: Good. You? Budi: Also good. Wati: What are you doing? Budi: I am reading a novel. Wati: What novel are you reading? Budi: I am reading a novel about everything about Islam! Wati: May I borrow it? Budi: Of course. Wati: It is already time for me to go home. Goodbye, and don't forget I'll borrow your book tomorrow Budi: Of course I won't forget. Goodbye too. Wati: Good bye. Selamat. The word "selamat" means "safe". So, "selamat pagi" literally means "good morning". The greetings in Indonesian are not quite the same as that of English. Below is the table of words with their meaning and the time you may want to use it: Unlike English, it is all right to greet people with "selamat malam" when meeting at night. To say "good bye", we can use "selamat tinggal". The word "selamat" also means "congratulations". Therefore, it is also used to congratulate other people. So, you can use the word "selamat" with the following words: So, "selamat ulang tahun" means "happy birthday". And so forth. Note that the word "jalan" means "street" or "to go", but when used in "selamat jalan", it means "bon voyage". Apa Kabar "(How Are You?)". The phrase "apa kabar" literally means "what (your) news". The word "apa" means "what" and the word "kabar" means "news". When translated, it means "how are you". To answer "apa kabar", we usually use "baik" or "baik-baik" to indicate that it's good. We can answer "biasa saja" (= "so so") or "kurang baik" (= "not good", literally = "less good"). In Malay, they use the spelling of "khabar" instead of "kabar", and thus the pronunciation is slightly different, with the Malay pronunciation using a hard H sound instead of regular K sound. Sapaan "(Salutations)". Notice that the dialog above uses "pak" and "bu", which mean "sir" and "ma'am" respectively. In Indonesian, you'll need to specify proper salutations in most cases when greeting people. This is because Indonesian people tend to be very polite. In formal situations or the work place, adults usually greet using "pak" or "bu". The word "Anda" (usually capitalised to show respect) is the general, relatively polite form of "you"; note that "bapak" and "ibu" could also be used, as well as casual forms such as "kamu". The article "pak" is shorthand for "bapak" (= mister or father), while "bu" is an abbreviation of "ibu" (= madam or mother). &lt;hr&gt; ^ Indonesian ^ | « Lesson 0: The Alphabet | Lesson 1: Greetings |Lesson 2: This, That » 

^ Indonesian ^ | « Lesson 1: Greetings | Lesson 2: This, That | Lesson 3: Pronouns » Arti Umum "(General Meaning)". The word "ini" means "this" or "these". The word "itu" means "that" or "those". Penggunaan Pertama: Ini buku "(First Usage: This [is a] book)". Contoh "(Example)": From the examples above, you can see that there is no explicit word for "to be" as in English. It is already implied by the structure or the vocabulary used in the sentence. In addition, when using a noun, the meaning can be either singular or plural. It depends on the context. When the context does not make it clear, the noun is assumed to be singular. The reader might ask: If "ini buku" means this is a book while when naively translated sounds like "This book." Then how to say "this book" in Indonesian? Answer: "buku ini". When learning Indonesian, or when Indonesian learning English, people usually realized that the order of words, especially adjective and noun, are usually swapped. Penggunaan Kedua: Buku ini merah "(Second Usage: This book [is] red)". Contoh "(Example)": Notice the difference between the first pattern and the second one: Bentuk Negatif "(Negative Form)". There are two words in Indonesian used to indicate negatives: In some cases, we can use either "bukan" or "tidak", but we'll discuss this later. Contoh "(Example)": Here, we use "bukan", because we're negating the nouns. However, look at the following example: This time, we are negating the adjective "right". (Benar = right / correct) Likewise, look at the following examples: Contoh "(Example)": In these examples, we use "tidak" because this time we're negating the adjectives. Bentuk Interogatif "(Interrogative Form)". The word "apa" (= "what") can be used to form questions using "ini" or "itu". Contoh "(Example)": You can use "apa ini" or "apa itu" to ask what something is called in Indonesian (by pointing to the object). ^ Indonesian ^ | « Lesson 1: Greetings | Lesson 2: This, That | Lesson 3: Pronouns » 

^ Indonesian ^ | « Lesson 2: This, That | Lesson 3: Pronouns |Lesson 4: Simple Sentences » General Pronoun. In Indonesian, both subjective and objective pronouns are the same. Possessive pronouns are slightly different in informal situations only. Below is the table: Note that there are two notions of "we" in Indonesian. If the opposite party is included, then we use "kita". Otherwise, we use "kami". Contoh "(Example)": Here, Budi speaks to Wati that this book is both Budi's and Wati's book. However, if Budi said, "Wati, ini buku kami."; it means that this book is Budi's (and probably his other friend's), but not Wati's. To refer a third person that has already died, He/She/It, when person we are speaking about already passed away and we want to refer to those people with respect, we use another word "almarhum" meaning "he" when the he is already died and "he" happens to be a respectable person. For female we use the word "almarhumah". This is a rare instance when Bahasa Indonesia is gender specific. This is a loan word from Arabic which literally means: "who was blessed by God." In this sense, this is actually a euphemism. So we can conclude that words 'almarhum' and 'almarhumah' are purposed to Muslims. The non-Muslim people usually change it with 'mendiang', which is simpler and can be used by both genders. Kata Kepunyaan "(Possessive Pronoun)". As you probably have noted in the previous example, the position of possessive pronoun is reversed in Indonesian. Contoh "(Example)": In these examples, notice that when using informal possessive of singular person, the suffix is put together with the noun. Although in spoken Indonesian it is acceptable to say "gelas dia" instead of "gelasnya" and "bolpen kamu" instead of "bolpenmu", it is incorrect to say "jeruk aku" to mean "my orange". Kata Kepunyaan #2 "(Possessive Pronoun Part 2)". In most languages there is the possibility of both adjective-like possessives "this is my book" and noun like adjectives "this book is mine" or "this is mine." Indonesian doesn't do this exactly like this but does have an equivalent: Notice that this kind of structure is "bridged" by the word ""punya"", which means "to have". You can then put the appropriate suffix or word to indicate the ownership. A synonym of ""punya" is "milik"". Hence you can change "punya" with "milik". Example: Contoh "(Examples)". Below are some examples that summarise what we've learnt: Note that at the third example we use "bukan" to deny the noun (i.e. the pen). At the fourth example, we use "tidak" to deny the ownership, which is considered as a verb (i.e. "doesn't have"). Note also how we can put names in ownership at example five and six with the word "punya". &lt;hr&gt; ^ Indonesian ^ | « Lesson 2: This, That | Lesson 3: Pronouns |Lesson 4: Simple Sentences » 

^ Indonesian ^ | « Lesson 3: Pronouns | Lesson 4: Simple Sentences |Lesson 5: Numbers » Tenses. Basic Indonesian word order is similar to English. Generally, sentences begin with a subject, followed by a verb (also called a predicate), and then an object. It's good news that Indonesian verbs don't change depending on tense. Indicating the past or future tense only requires inserting words that indicate the time, in a very regular system. Contoh (Example): As you noticed in the examples above, the word "telah" or "sudah" indicate completed actions, the word "akan" indicate future actions, and the word "sedang" indicate actions in progress. The main verb (i.e. "makan" = to eat) is left unchanged. Side note: In English, both cooked and uncooked rice are referred as rice alone. In Indonesian, uncooked rice is called beras while cooked rice is called nasi. Note also that Indonesian has no notions of imperfect tense. Instead, Indonesian uses duration words, such as "selama" (literally means during or as long as): Contoh (Example): Don't worry about these "tenses" yet. We'll study it in greater detail later. Adding Emphasis. The Indonesian language is very expressive. While basic Indonesian word order matches English, you can scramble up the sentence structure, and the sentence will still have the same underlying meaning. In this regard, Indonesian is somewhat like Latin or Japanese, but without the cases or the particles. Usually, when a word other than the subject is put at the beginning of a sentence, it becomes the emphasis of the sentence. This is broadly used in spoken Indonesian. Contoh (Example): Note also that this way of providing emphasis can occur with any "tense" using the same pattern. Don't be intimidated by this variability in word order. You can always form simple sentences as: Subject + Verb + Object. &lt;hr&gt; ^ Indonesian ^ | « Lesson 3: Pronouns | Lesson 4: Simple Sentences |Lesson 5: Numbers » 

^ Indonesian ^ | « Why learn Indonesian? | How to use this Indonesian Wikibook |Lesson 0: The Alphabet » Cara Pemakaian "(Usage)". I plan to structure this book so that for each lesson you will have the corresponding exercise. So, if you have finished lesson 1, do the exercise 1 before proceeding to lesson 2, and so forth. This is so that you fully understand the Indonesian. I understand that it is pretty hard to illustrate simple concepts of Indonesian without knowing anything about it. So, my approach is to illustrate the most important basic components of Indonesian in the basic section. My suggestion is to skim the basic section and grasp a little bit about it and step up to the beginner section so that you can see these components in action. ^ Indonesian ^ | « Why study Indonesian? | How to use this Indonesian Wikibook |Lesson 0: The Alphabet » 

^ Indonesian ^ | « Lesson 4: Simple Sentences | Lesson 5: Numbers |Lesson 6: Particles » Angka Numerals "(Numerals)". Contoh "(Examples)": Urutan "(Series)". As you see, in order to form ordinal numbers in series, you only need to attach the prefix "ke-", then the number itself. Except for pertama ("1st"), you can spell out the number or just write the digits (like the ones in parentheses). Contoh (Examples): Nilai Tempat "(Value Place)". Contoh "(Examples)": Berapa "(How much)". To ask about quantities, we use the word "berapa". For example: "Berapa harga X" is to ask "How much does X cost". Contoh (Examples): The word "berapa" can directly replace the quantity in question. However, you must put a measure word right after it. Contoh (Examples): The most common measure word is "buah", which indicates quantity. Coincidentally, this word can also mean fruit (i.e. apples, oranges, etc). Other common measure words are "ekor" (for animals; ekor = tail) and "orang" (for people, orang = person). Interestingly, the word biji (lit= fruit) is also sometime use to replace buah. The examples below illustrate using measure words: Contoh (Examples): Contoh Lain (Other Examples): Angka Lebih Dari 999 "(Numbers higher than 999)". In English, to make the number readable, we used to add a comma every three digits. For example, One million and two hundred and fifty-five thousands and three hundred and sixty-four can be written as 1,255,364. However in Indonesian, we insert a dot instead of comma. Hence the same number (Indonesia: "Satu juta dua ratus lima puluh lima ribu tiga ratus enam puluh empat") should be written as 1.255.364 . For curious reader, Indonesians use comma (",") to indicate decimals. "Online Tools". &lt;hr&gt; ^ Indonesian ^ | « Lesson 4: Simple Sentences | Lesson 5: Numbers |Lesson 6: Particles » 

^ Indonesian ^ Kata Pembilang - "Quantification Words". Indonesian language has no distinction between uncountable and countable nouns. If you want to say that there are more than one of a thing, as a rule of thumb you can repeat the noun using a dash (-) to make it plural. Alternatively you can also put the total number of the thing in front of the noun. There is no rule in Indonesian language regarding which nouns should/can be repeated, which nouns should not/cannot, but if we already put some quantification words, we can't do the repetition noun. For example: It is wrong to say banyak buku-buku - "lots of books", instead use either the phrase banyak buku or buku-buku These words are essential to specify indefinite plurals: Kata Penunjuk Jumlah - Measure Words. An example of a measure words is the word "pieces" in "two pieces of paper". You have seen the word buah in lesson 5 to indicate the total of a quantity. The word buah can be considered as the "universal" measure word because it can fit in so many situations. However, you may want to fine tune it to convey a more precise meaning. For example: in the English language it's more proper to say "two rolls of bread" or "two slices of bread" rather than "two pieces of bread". This is because the measure word "roll" is more fitting than the word "piece" in this situation. The same thing applies to the Indonesian language. Below are other common measure words, followed by the noun it measures. &lt;hr&gt; ^ Indonesian ^ 

^ Indonesian ^ | « Lesson 5: Numbers | Lesson 6: Particles |Lesson 7: Introducing Yourself » Di "(In/on/at [place])". To indicate a place, we use the particle di. It can mean in, on, or at. Contoh (Examples): Note the word "ada", which means "to exist". It is placed right before the di particle to indicate existence. Note also that the particle di doesn't convey any further detail on how the object is being placed, whether it's in front, inside, etc. To put additional detail, we put a location word after the particle di. Contoh (Examples): The word "atas" means "top" or "above" and "dalam" means "inside". Below is the list of location words you may use: In Indonesian, to change the position is quite easy. Like Buku saya ada di atas meja, if you want to change it becomes beside or behind, you just need to change them into Buku saya ada di sebelah/samping meja and Buku saya ada di belakang meja. Be careful to differentiate between di as prefix (awalan) and di as showing the place (kata depan). Most Indonesian natives forget about this and mistakes are common. If the word di is followed by a verb, it is a prefix. Example Indonesian also commonly (and mistakenly) use di with time, for example "di waktu sedih" (during sad times). Pada is the correct proposition for time. Hence, it should be "pada waktu sedih". Pada "(at [person])". Although in spoken Indonesian it is acceptable to use "di" to indicate the existence of a noun at someone, this is unusual. For example, it is not correct to say: "Bukumu ada di Budi" to mean "Your book is at Budi". Rather, you should use the particle pada: This sounds awkward to translate literally. Usually in English, people would say "Budi has your book". This preposition is also used for time, for example Pada pukul enam pagi. (At 6 am) Ke "(to [a place])". The particle ke is to indicate the notion of to a place. It is often coupled with the word "pergi", which means "to go". Contoh (Example): In spoken Indonesian, people often omit pergi when the context is clear. So, you'll often hear Ibu ke pasar to mean Mother [goes] to [the] market. Kepada "(to [a person])". Some verbs in English like "to send", "to give" and so on need the particle "to", followed by a person. For example: "I give the book to you". In Indonesian, for this notion of to, you cannot use the particle ke. Rather, you'll use the particle kepada. Contoh (Example): Memberikan = to give Mengirimkan = to send Surat = letter Certainly, in spoken Indonesian, people may violate this rule and use ke instead of the proper kepada. Dari "(From)". The particle dari is almost synonymous with from in English. It is to indicate the origin of something. Contoh "(Example)": Not only that dari explains the place origin, but also explains the origin of things. For example: Cincin = ring Terbuat = is made Emas = gold Untuk "(For)". The particle untuk is almost synonymous with for. For example: It is also used to explain the usage of a thing: &lt;hr&gt; ^ Indonesian ^ | « Lesson 5: Numbers | Lesson 6: Particles |Lesson 7: Introducing Yourself » 

Puzzles | Logic puzzles | Knights, Knaves &amp; Spies II | Solution C should answer "No". Reasoning. B is either a knave or a spy. If B is a spy, then A is truthful and is therefore the knight.  A - Knight  B - Spy  C - Knave On the other hand, if B is the knave, there are two possibilities:  A - Spy  B - Knave  C - Knight or  A - Knight  B - Knave  C - Spy If C is either the knave or the knight, his answer to the question will be "No", and so the judge will not be able to draw a conclusion. On the other hand, C can answer "Yes" only if he is the spy. 

^ Indonesian ^ Kata Tanya "(Question Words)". The following is a list of Indonesian question words with their use. Yang mana - "Which [one]". &lt;hr&gt; ^ Indonesian ^ 

Human civilization in the Americas probably began in the last ice age, when prehistoric hunters crossed a land bridge between the Asian and North American continents. Civilizations in North America, Central America, and South America had different levels of complexity, technology, and cohesiveness. Some of the most powerful and organized societies occurred in South and Central America. These cultures developed writing, allowing them to spread and dominate. They created some of the largest cities in the Ancient world. North American cultures were more fragmented and less unified. The tribe was often the major social unit, with exchanges between tribes creating similar societies over vast distances. Tribal dwellings as large as European towns flourished in the rugged desert of southwestern North America. European-descended historians have difficulty referring to these cultures as a whole, as the native people did not have a unified name for themselves. At first, Europeans called natives "Indians". This term came from the belief by Christopher Columbus that he had discovered a new passage to India. Despite Amerigo Vespucci ascertaining that the Americas were not actually India, Indian continued to be used as the 'de facto' name for native inhabitants until around 1960. Starting in the 1960s, the term "Native American" was used. Yet this term may be too problematic: The name America derives from Amerigo Vespucci, an Italian who had little to do with the native people. There is also "American Indian". This is too general a term for a group having little in common other than skin tone and non-European language. In Canada, the term "First People" is used. All these terms for the native people of America show just how diverse Pre-Columbian America was and the disagreement continues among scholars today about this period. Early Inhabitants of the Americas. Bering Land Bridge. American history does not begin with Columbus's 1492 arrival. The Americas were settled long before the first European arrived. Civilization began during the last ice age, some 15 to 40 thousand years ago. Huge ice sheets covered the north, so sea levels were much lower, creating a land bridge between Asia and North America. This was the Bering land bridge, a gap in two large ice sheets creating a connection from lands near present day Alaska, through Alberta, and into the continental United States. Nomadic Asians following herds of wild game traveled into the continental United States. A characteristic arrow point was found and first described near the present day town of Clovis, NM. Specialized tools and common burial practices are seen in many archaeological sites through North America and into South America. Clovis People. The Clovis people are one of North America's earliest civilizations. It is not clear if the finds represent one unified tribe, or many tribes with a common technology and belief. Their trek across 2000 rugged miles is one of the great feats of pre-history. Their culture disappears dramatically from the archaeological record 12,900 years ago, with widespread speculation about what caused their disappearance. Theories range from the extinction of the mammoth, to sudden environmental changes caused by a comet, to flooding caused by the break of a massive freshwater lake, Lake Agassiz. There is controversy about Pre-Clovis settlement of North and South America. Comparisons of culture and linguistics offer evidence of the influence of early America by several different contemporary cultures. Some genetic and time-dating studies point to the possibility that ancient Americans came from other places and arrived earlier than at the Clovis sites in North America. Perhaps some ancient settlers to the hemisphere traveled by boat along the seashore, or arrived by boats from the Polynesian islands. As time went on, many of these first settlers settled down into agricultural societies, complete with domesticated animals. Groups of people formed stable tribes and developed distinct languages of their own, to the point that more distant relatives could no longer understand them. Comparative linguistics -- the study of languages of different tribes -- shows fascinating diversity, with similarities between tribes hundreds of miles apart, yet startling differences with neighboring groups. At times, tribes would gain regional importance and dominate large areas of America. Empires rose across the Americas that rivaled the greatest ones in Europe. For their time, some of these empires were highly advanced. Early Empires of Mesoamerica. Meso-American civilizations are among some of the most powerful and advanced civilizations of the ancient world. Reading and writing were widespread throughout Meso-America, and these civilizations achieved impressive political, artistic, scientific, agricultural, and architectural accomplishments. Many of these civilizations gathered the political and technological resources to build some of the largest, most ornate, and highly populated cities in the ancient world. Maya. The aboriginal Americans settled in the Yucatan peninsulas of present-day Mexico around 10,000 BCE. By 2000 BCE, the Mayan culture had evolved into a complex civilization. The Mayans developed a strong political, artistic and religious identity among the highly populated Yucatan lowlands. The classic period (250-900 CE) witnessed a rapid growth of the Mayan culture and it gained dominance within the region and influence throughout present-day Mexico. Large, independent city-states were founded and became the political, religious, and cultural centers for the Mayan people. Mayan society was unified not by politics, but by their complex and highly-developed religion. Mayan religion was astrologically based, and supported by careful observations of the sky. The Mayans had a strong grasp of astronomy that rivaled, and, in many ways, exceeded that of concurrent European societies. They developed a very sophisticated system for measuring time, and had a great awareness of the movements in the nighttime sky. Particular significance was attached to the planet Venus, which was particularly bright and appeared in both the late evening and early morning sky. Mayan art is also considered one of the most sophisticated and beautiful of the ancient New World. The Mayan culture saw a decline during the 8th and 9th century. Although its causes are still the subject of intense scientific speculation, archaeologists see a definite cessation of inscriptions and architectural construction. The Mayan culture continued as a regional power until its discovery by Spanish conquistadors. In fact, an independent, non-centralized government allowed the Mayans to strongly resist the Spanish conquest of present-day Mexico. Mayan culture is preserved today throughout the Yucatan, although many of the inscriptions have been lost. Aztec. The Aztec culture began with the migration of the Mexica people to present-day central Mexico. The leaders of this group of people created an alliance with the dominant tribes forming the Aztec triple alliance, and created an empire that influenced much of present-day Mexico. The Aztec confederacy began a campaign of conquest and assimilation. Outlying lands were inducted into the empire and became part of the complex Aztec society. Local leaders could gain prestige by adopting and adding to the culture of the Aztec civilization. The Aztecs, in turn, adopted cultural, artistic, and astronomical innovations from its conquered people. The heart of Aztec power was economic unity. Conquered lands paid tribute to the capital city Tenochtitlan, the present-day site of Mexico City. Rich in tribute, this capital grew in influence, size, and population. When the Spanish arrived in 1521, it was the fourth largest city in the world (including the once independent city Tlatelolco, which was by then a residential suburb) with an estimated population of 212,500 people. It contained the massive Temple de Mayo (a twin-towered pyramid 197 feet tall), 45 public buildings, a palace, two zoos, a botanical garden, and many houses. Surrounding the city and floating on the shallow flats of Lake Texcoco were enormous "chinampas" -- floating garden beds that fed the many thousands of residents of Tenochtitlan. While many Meso-American civilizations practiced human sacrifice, none performed it to the scale of the Aztecs. To the Aztecs, human sacrifice was a necessary appeasement to the gods. According to their own records, one of the largest slaughters ever performed happened when the great pyramid of Tenochtitlan was reconsecrated in 1487. The Aztecs reported that they had sacrificed 84,400 prisoners over the course of four days. With their arrival at Tenochtitlan, the Spanish would be the downfall of Aztec culture. Although shocked and impressed by the scale of Tenochtitlan, the display of massive human sacrifice offended European sensitivity, and the abundant displays of gold and silver inflamed their greed. The Spanish killed the reigning ruler, Montezuma in June 1520 and lay siege to the city, destroying it in 1521, aided by their alliance with a competing tribe, the Tlaxcala. Inca. With the ascension of Manco Capac to emperor of a tribe in the Cuzco area of what is modern-day Peru around 1200 BCE, the Incan civilization emerged as the largest Pre-Columbian empire in the Americas. Religion was significant in Inca life. The royal family were believed to be descendants of the Inca Sun God. Thus, the emperor had absolute authority, checked only by tradition. Under the emperors, a complex political structure was apparent. The Incan emperor, regional and village leaders, and others were part of an enormous bureaucracy. For every ten people, there was, on average, one official. The organization of the Empire also included a complex transportation infrastructure. To communicate across the entire empire, runners ran from village to village, relaying royal messages. In 1438, the ambitious Pachacuti, likely the greatest of the Incan emperors, came to the throne. Pachacuti rebuilt much of the capital city, Cuzco, and the Temple of the Sun. The success of Pachacuti was based upon his brilliant talent for military command (he is sometimes referred as the "Napolean of the Andes") and an amazing political campaign of integration. Leaders of regions that he wanted to conquer were bribed with luxury goods and enticed by promises of privilege and importance. As well, the Incans had developed a prestigious educational system which, not incidentally, just happened to extol the benefits of Incan civilization. Thus, much of the expansion throughout South America was peaceful. At its height of power in the late 15th century, Incan civilization had conquered a vast patchwork of languages, people and cultures from present-day Ecuador, along the whole length of South America, to present-day Argentina.Cuzco, the capital city, was said by the Spanish to be "as fine as any city in Spain". Perhaps the most impressive city of the Incan empire, though, was not its capital, Cuzco, but the city Machu Picchu. This mountain retreat was built high in the Andes and is sometimes called the "Lost City of the Incas." It was intended as a mountain retreat for the leaders of the Incan empire and demonstrates great artistry -- the abundant dry stone walls were entirely built without mortar, and the blocks were cut so carefully that one can't insert a knife-blade between them. The Spanish discovered the Inca during a civil war of succession and enjoyed great military superiority over the slow siege warfare that the Incan empire had employed against its enemies. Fueled by greed at the opportunity to plunder another rich civilization, they conquered and executed the Incan emperor. The Incan empire fell quickly in 1533, but a small resistance force fled to the mountains, waging a guerrilla war of resistance for another 39 years. Meso-American Empires. The Meso-American Empires were undoubtedly the most powerful and unified civilizations in the new world. Writings were common in Meso-America and allowed these cultures to spread in power and influence with far more ease than their counterparts in north America. Each of these civilizations built impressive urban areas and had a complex culture. They were as 'civilized' as the Spanish who conquered them in the 15th and 16th centuries. Early Empires of the Southwest. Native Americans adapted the arid desert southwest. A period of relatively wet conditions saw many cultures in the area flourish. Extensive irrigation was developed that were among the biggest of the ancient world. Elaborate adobe and sandstone buildings were constructed. Highly ornamental and artistic pottery was created. The unusual weather conditions could not continue forever, though, and gave way, in time, to the more common drought of the area. These dry conditions necessitated a more minimal way of life and, eventually, the elaborate accomplishments of these cultures were abandoned. Ancestral Puebloans. One prominent group were the Ancestral Puebloans, who lived in the present day Northeastern Arizona and surrounding areas. The geography of this area is that of a flat arid, desert plain, surrounded by small areas of high plateau, called mesas. Softer rock layers within the mesas eroded to form steep canyons and overhangs along their slopes. The Ancestral Puebloans culture used these cave-like overhangs in the side of steep mesas as shelter from the brief, fierce southwestern storms. They also found natural seeps and diverted small streams of snow melt into small plots of maize, squash and beans. Small seasonal rivers formed beds of natural clays and dried mud. The Ancestral Puebloans used hardened dry mud, called adobe, along with sandstone, to form intricate buildings that were sometimes found high in the natural overhangs of the mesas. The Ancestral Puebloans were also skilled at forming the natural clays into pottery. Between 900 - 1130 CE a period of relatively wet conditions allowed the Ancestral Puebloans to flourish. Traditional architecture was perfected, pottery became intricate and artistic, turkeys were domesticated, and trade over long distances influenced the entire region. Following this golden period was the 300 year drought called Great Drought. The Ancestral Puebloan culture was stressed and erupted into warfare. Scientists once believed the entire people vanished, possibly moving great distances to avoid the arid desert. New research suggests that the Ancestral Puebloans dispersed; abandoning the intricate buildings and moving towards smaller settlements to utilize the limited water that existed. Hohokam. Bordering the Ancestral Puebloan culture in the north, a separate civilization emerged in southern Arizona, called the Hohokam. While many native Americans in the southwest used water irrigation on a limited scale, it was the Hohokam culture that perfected the technology (all without the benefit of modern powered excavating tools). The ability to divert water into small agricultural plots meant that the Hohokam could live in large agricultural communities of relatively high population density. This was particularly true in the Gila River valley, where the Gila River was diverted in many places to irrigate large fertile plains and numerous compact towns. The bigger towns had a 'Great House' at their centers, which was a large Adobe/stone structure. Some of these structures were four stories in size and probably were used by the managerial or religious elites. Smaller excavation or pits were enclosed by adobe walls and used as primary residences. Smaller pit rooms and pits were used for many different functions. The successful use of irrigation is evident in the extensive Casa Grande village. Situated between two primary canals, the Casa Grande site has been the focus of nearly 9 decades of archaeological work. The original town was built around a Great House and incorporated open courtyards and circular plazas. By the 10th century neighboring settlements had been built to accommodate a large, highly developed region. The scale of this community can be seen in the results of one excavation of part of it in 1997. The project identified 247 pit houses, 27 pit rooms, 866 pits, 11 small canals, a ball court, and portions of four adobe walled compounds. The Hohokam culture disintegrated when they had difficulty maintaining the canals in the dry conditions of the drought. Small blockages or collapses of the canal would choke the intricate irrigation networks. Large towns and extensive irrigation canals were abandoned. The people gave up their cultural way of life and dispersed into neighboring tribes. Early Empires of the Mississippi. Native Americans in the Eastern Continental United States developed mound-building cultures early in North American History. Groups of native Americans became more stratified as time went on and developed into tribes. These tribes participated in long networks of trade and cultural exchanges. The importance of trade routes developed urban cities of great influence. The mound-building people were one of the earliest civilization to emerge in North America. Beginning around 1000 BCE cultures developed that used mounds for religious and burial purposes. These mound-building people are categorized by a series of cultures that describe distinctive artwork and artifacts found in large areas throughout the present-day eastern United States. The burial mound was the principle characteristic of all of these societies. These large structures were built by piling baskets of carefully selected earth into a mound. Mounds were pyramid shaped with truncated tops. Sometimes small buildings were built on top of them. Some of these mounds were quite large. The Grave Creek Mound, in the panhandle of present day West Virginia, is nearly 70 feet tall and 300 feet in diameter. Other mounds have even been shown to be oriented in a way that allows for astronomical alignments such as solstices and equinoxes. Mound building cultures spread out in size and importance. The first culture, the Adena, lived in present-day Southern Ohio and the surrounding areas. The succeeding cultures united to create an impressive trade system that allowed each culture to influence the other. The Hopewell exchange included groups of people throughout the continental Eastern United States. There began to be considerable social stratification within these people. This organization predates the emergence of the tribe as a socio-political group of people that would dominate later eastern and western native American civilization. The climax of this civilization was the Mississippian culture. The mound-building cultures had progressed to social complexity comparable to Post-Roman, Pre-Consolidation Tribal England. Mounds became numerous and some settlements had large complexes of them. Structures were frequently built on top of the mounds. Institutional social inequalities existed, such as slavery and human sacrifice. Cahokia, near the important trade routes of the Mississippi and Missouri rivers, became an influential and highly developed community. Extensive trade networks extended from the Great Lakes to the Gulf of Mexico. Cahokia was one of the great centers of Mississippian culture and its largest settlement of Mississippi. The focal point of the settlement was the ceremonial mound called Monk's Mound. Monk's Mound was the largest mound ever constructed by mound-building people and was nearly 100 feet tall and 900 feet long. Excavation on the top of Monk's Mound has revealed evidence of a large building - perhaps a temple - that could be seen throughout the city. The city was spread out on a large plain south of Monks mound. The city proper contained 120 mounds at varying distances from the city center. The mounds were divided into several different types, each of which may have had its own meaning and function. A circle of posts immediate to Monk's Mound marked a great variety of astronomical alignments. The city was surrounded by a series of watchtowers and occupied a diamond shape pattern that was nearly 5 miles across. At its best, the city may have contained as many as 40,000 people, making it the largest in North America. It is likely the Mississippian culture was dispersed by the onslaught of viral diseases, such as smallpox, which were brought by European explorers. Urban areas were particularly vulnerable to these diseases, and Cahokia was abandoned in the 1500's. The dispersal of tribes made it impractical to build or maintain mounds and many were found abandoned by European explorers. Contact with European Culture. Epidemics. European contact brought immediate changes in many tribes of North America. One of the most significant changes to all Indian tribes was the introduction of viral diseases and epidemics. Smallpox was probably the single biggest scourge to hit North America. Infected contagious Indians spread the plague far inland almost immediately after early encounters with European settlers. It is estimated that around 90% of all Native Americans died from diseases soon after first contact. The effects traumatized many powerful and important cultures. Urban areas were particularly vulnerable and Native American culture adapted by becoming more isolated, less unified, and with a renewed round of inter-tribal warfare as tribes seized the opportunity to gain resources once owned by rivals. Columbian Exchange. On the other hand, Europeans brought invasive plants and animals. The horse was re-introduced to America (as original paleo-American populations of wild horses from the Bering land bridge were extinct) and quickly adapted to free range on the sprawling great plains. Tribes of nomadic Native Americans were quick to see the horse's value as an increase in their mobility; allowing them to better adapt to changing conditions and as a valuable asset in warfare. Along with Europeans bringing plants and animals, the Europeans were able to take several plants such as corn, potatoes, and tomatoes back to their native countries. Review Questions. 1. Give two names for the indigenous peoples living in America, and name the circumstances behind each name. 2. What evidence do we have for the Inca, Mayan, and Aztec cultures? 3. What in the climate contributed to the rise and fall of the indigenous peoples of South-West North America? 



Puzzles | Logic puzzles | Lying about your Age | Solution Annie - 30&lt;br&gt; Betty - 51&lt;br&gt; Carrie - 55&lt;br&gt; Darla - 46&lt;br&gt; Eve - 37 Reasoning. Symbolic. Let the ages and names of Annie, Betty, Carrie, Darla and Eve be A, B, C, D and E. C says to A, that C = A + 10. If C were younger than A, that would be lying, so C must be older than A. (But still lying.) We have A &lt; C. C says to A, that B &lt; D. As C &gt; A, C is lying, so B &gt; D. We have A &lt; C, D &lt; B. D says to B, that D = E + 9. As D &lt; B, D is telling the truth, so D &gt; E. We have A &lt; C, E &lt; D &lt; B, D = E + 9. E says to B, that E = A + 7. As E &lt; B, E is telling the truth, so E &gt; A. We have A &lt; C, A &lt; E &lt; D &lt; B, D = E + 9, E = A + 7. Since D = E + 9 and E = A + 7, D = A + 7 + 9 = A + 16. We have A &lt; C, A &lt; E &lt; D &lt; B, D = E + 9 = A + 16, E = A + 7. B says to C, that E &lt; C. If B &gt; C then B would be lying, so then E &gt; C, and then A &lt; C &lt; E &lt; D &lt; B. However, C says to D, that C = D ± 6; since C &lt; D, this gives C = D - 6. However, we have E = D - 9, which would make E &lt; C, giving a contradiction. The assumption that B &gt; C is therefore false, so B &lt; C. We have A &lt; E &lt; D &lt; B &lt; C, D = E + 9 = A + 16, E = A + 7. A says to B, that B = (17/10)A. As A &lt; B, A is telling the truth. We have A &lt; E &lt; D &lt; B &lt; C, B = (17/10)A, D = E + 9 = A + 16, E = A + 7. B says to C, that |C - D| = |D - E| → |C - D| = 9. As B &lt; C, B is telling the truth, so C = D + 9. As D = A + 16, C = A + 16 + 9 → C = A + 25. We have A &lt; E &lt; D &lt; B &lt; C, B = (17/10)A, C = A + 25, D = A + 16, E = A + 7. Using D &lt; B &lt; C, we have A + 16 &lt; (17/10)A &lt; A + 25 → 16 &lt; (7/10)A &lt; 25 → 160/7 &lt; A &lt; 250/7 → 22 + 6/7 &lt; A &lt; 35 + 5/7. Since B and A must both be whole numbers, and B = (17/10)A → B - A = (7/10)A, (7/10)A must be a whole number. Hence A must be divisible by 10. The only whole number fitting 22 + 6/7 &lt; A &lt; 35 + 5/7 is A = 30. We have A = 30, B = (17/10)A, C = A + 25, D = A + 16, E = A + 7. Hence A = 30, B = 51, C = 55, D = 46, E = 37. Verbal. Carrie tells Annie she's older than her by 10 years. If Carrie is younger, she's lying, and that's impossible, so Carrie must be older than Annie, just not by 10 years. FACT: Carrie is older than Annie (but not by 10 years). Carrie also lies to (younger) Annie that Betty is younger than Darla. FACT: Darla is younger than Betty. Darla tells the truth to (older) Betty that she's 9 years older than Eve. FACT: Darla is 9 years older than Eve. Eve tells the truth to (older) Betty that she's 7 years older than Annie. FACT: Eve is 7 years older than Annie. Annie tells the truth to (older) Betty that Betty's age is 70% greater than her own. For Betty's age to be a whole number, Annie's age must be a multiple of 10. Since Betty is older than Darla, and Darla is 7 + 9 = 16 years older than Annie, that means Betty has to be more than 16 years older than Annie. The lowest multiple of 7 greater than 16 is 21. FACT: Annie is at least 30 years old (and definitely a multiple of 10). At this point, Betty appears to be the oldest, lying lady. Let's assume that, and see if it works. In that case, Carrie is lying to Darla that the difference in their ages is 6 years, but Betty tells the truth to (older) Carrie that the difference between Carrie's age and Darla's is the same as the difference between Darla's and Eve's, namely, 9 years. Let's test this scenario, assuming Annie's age is 30. Then we get, from youngest to oldest: TESTING: Annie = 30, Eve = 37, Darla = 46, Betty = 51, Carrie = 55 Checking all statements and the age relations shows that this is an answer. Is this the only answer? If Annie's age was 40, then Betty's age would be 68, and Carrie's age would be 65, so Carrie would not be the oldest, and that would be a fatal flaw. If Annie is older than 30, Betty is older than Carrie, and Carrie is not the oldest. Hence, it must have been the only answer. 

Puzzles | Logic puzzles | Getting out of Prison | Solution "Out of the answers 'left' and 'right', which door would the other guard say doesn't lead to freedom, if I asked her?" - The answer will be the door that leads to freedom. Both guards will answer the same. OR similarly "Would the other guard tell me to go through your door to get to freedom?" If either guard says "Yes," go through the other guard's door. If either says "No," go through her door.. OR alternatively Trick question: "If I asked you whether the door you guard is the execution room, would you say yes?" If the guard you asked was guarding the freedom room, the lying guard would have said Yes to the simple question "Is this the execution room?", but when asked whether he would say yes to this simple question, he has to lie and say No to the Trick question. If the guard you asked was guarding the freedom room, the true guard would have said No to the simple question "Is this the execution room?", and when asked whether he would say yes to this simple question, he will truthfully say No to the Trick question. Thus if it is the freedom door, either guard will give the answer No to the Trick question. If the guard you asked was guarding the execution room, the lying guard would have said No to the simple question "Is this the execution room?", but when asked whether he would say yes to this simple question, he has to lie and say Yes to the Trick question. If the guard you asked was guarding the execution room, the true guard would have said Yes to the simple question "Is this the execution room?", and when asked whether he would say yes to this simple question, he will truthfully say Yes to the Trick question. Thus if it is the execution door, either guard will give the answer Yes to the Trick question. Therefore: If the answer to the trick question by the guard you ask is No - then that guard is guarding the freedom door. If the answer to the trick question by the guard you ask is Yes - then that guard is guarding the execution door. You can ask simply to the either of the guard "Is this my execution day?" If he answers "no" then you probably should take the other way out.Or if he says yes he is telling the truth and you should take the path guarded by him. 

Puzzles | Logic puzzles | Hats | Solution The first man saw at most 1 blue hat in front of him, for if he saw 2 he would have said "red" (there are only 2 blue hats). The second man saw a red hat in front of him, for if he saw a blue hat he would have said "red" (the first man saw at most 1 blue hat). That's how the last man now knows the 2nd man saw a red hat on him. 

The man picks up a piece of wood and lights it from the fire on the west end of the island. He then carries it over to set fire to the middle of the island. Fortunately, the wind is hurricane strength, so the fire from the middle only burns the east side of the island. The man then goes to the east side of the island, where the fire from the west side can't spread, because the east side has already burned down. The man survives the fire, but dies of starvation, with all the food in the forest burnt. 

 *---*---*---  | \ /  * * *  | X  * * *  | / 

Puzzles | Geometric puzzles | Rubber Band | Solution 1 Stretch the rubber band, and draw a line along it on the plane. By stretching the rubber band, while holding the loose end to the nail, you can measure 4cm. By folding it again, you can measure 2cm. Similarly, you can measure 1cm and then 0.5cm. Mark the 0.5cm on the line, and using the last 0,5cm of the rubber band, mark the line you drew at 0.5cm intervals, so you can use it as a ruler. Draw two circles around the nail, with radius 1.5cm and 2cm. Using 6cm of the rubber band, make a triangle, with a corner touching each circle, and the last corner at the nail. You now have a right angle at the nail. 

&lt; Back to Problem Some answers are formula_1, formula_2, formula_3 and formula_4. 

The answer is 2, 2 and 9. The product of their ages is 36, leaving the possibilities Knowing the sum of their ages is not sufficient to determine their ages, eliminating {1,1,36} (38), {1,2,18} (21), {1,3,12} (16), {1,4,9} (14), {2,3,6} (11) and {3,3,4} (10), leaving only Knowing that there exists an eldest daughter eliminates {1,6,6}, leaving only {2,2,9}. How does knowing that there exists an eldest daughter eliminate {1,6,6}? One daughter can be older than another but both have the same age in years. Hey, questioner, we're going by the same number of years. We are taking that both daughters in the thrown out {1, 6, 6} choice are EXACTLY 6-years-old. 

Puzzles | Arithmetical puzzles | Digits of the Square | Solution Looking at the last digit, the last digit must be either 0, 1, 5 or 6. Then looking at the last two digits, the last two digits must be either 00, 01, 25 or 76. Then looking at the last three digits, the last three digits must be either 000, 001, 625 or 376. Then looking at the last four digits, the last four digits must be either 0000, 0001, 0625 or 9376. Out of those, only 9376 is a 4 digit number. 

Puzzles |Decision puzzles|12 Coins|Solution You need 3 weighings. There are 12 possible places for the fake coin to be and it can either be heavy or light, leaving 24 possibilities. Each weighing can give 3 results. (One side heavier, other side heavier, or both sides equal.) 2 weighings can at most only differentiate between 9 possibilities, therefore 2 or less weighings is "not" sufficient. Label the coins with the symbols αβγδABCD$£€¤. Start at state S00.  State | Weighings | Pile L | Pile R | if L&lt;R | if L=R | if L&gt;R | Candidates in  | left | | | | | | αβγδABCD$£€¤  S00 | 3 | αβγδ | ABCD | goto S01 | goto S02 | goto S03 | ????????????  S01 | 2 | αβC | Bγδ | goto S04 | goto S05 | goto S06 | LLLLHHHH  S02 | 2 | αβγ | $£€ | goto S07 | goto S08 | goto S09 | ????  S03 | 2 | αβC | Bγδ | goto S10 | goto S11 | goto S12 | HHHHLLLL  S04 | 1 | α | β | α light | B heavy | β light | LL H  S05 | 1 | A | D | D heavy | error | A heavy | H H  S06 | 1 | γ | δ | γ light | C heavy | δ light | LL H  S07 | 1 | $ | £ | £ heavy | € heavy | $ heavy | HHH  S08 | 1 | α | ¤ | ¤ light | error | ¤ heavy | ?  S09 | 1 | $ | £ | $ light | € light | £ light | LLL  S10 | 1 | γ | δ | δ heavy | C light | γ heavy | HH L  S11 | 1 | A | D | A light | error | D light | L L  S12 | 1 | α | β | β heavy | B light | α heavy | HH L Reasoning. Since one of the coin is fake, and it can be lighter or heavier, that gives 2 X 12 = 24 possible answers. Also notice that each weighing can have 3 possible outcomes (left lighter, right lighter, equal), and with 3 weighting there are 33 = 27 possible outcome. Because 27 &gt; 24, it is possible (although not certainly) to have a solution. First weighing. Suppose we put x coins on both side of the weigh, we want to ensure that no matter what the outcome is, the remaining case still satisfy the condition of possible outcomes &gt; possible answers, otherwise the plan fails. In case the result is equal, only the coins not on weigh can be faked, since we don't know yet whether it is lighter or heavier, that leaves 2(12-2x) possible answers. With 2 weighing lefts, there are 32 = 9 possible outcomes. Hence 2(12-2x)&lt;9. When one side is lighter, it can be that one of the coins on the lighter side is lighter, or one of the coins on the heavier side is heavier, which leaves 2x possible answers. Hence we have 2x&lt;9. Solving the equation, we get x=4. 

Puzzles | Puzzles/Decision puzzles | Monty Hall | Solution There is a prior 1/3 chance of the door you picked being the door with the prize, and a 2/3 chance of a different door being the door with the prize. Therefore, if you want the prize, it is better to switch, and if you want a nice goat, it is better to stick with the door you chose. Reasoning. There are three doors: the door you first chose, the door that Monty opened, and the third door. Because there are two doors left, one with the prize, then one might want to believe that there is a 1/2 chance of obtaining the prize either way. However, before Monty opened any door, there was a 1/3 chance that the door contained the prize. When Monty opens a different door, you learn nothing about the door that you first chose. Thus, the chance of the door you first chose containing the prize remains 1/3, not 1/2. However, when Monty opens a door with the goat, the probability of that door containing the prize drops to 0. As the probabilities must sum to 1, and the probability of a prize at the door you first chose remains 1/3, the probability of a prize at the third door rises to 2/3. This rise in probability results only because you learned that Monty decided not to open that particular door. Calculations. Let formula_1 be the door you first picked; let formula_2 be the door which Monty opened; let formula_3 be the third and remaining door. formula_4 is a car or other prize and formula_5 is any goat. Suppose you choose the third door formula_3, then what is the probability formula_7 that this door contains the prize? formula_8 Line-by-line, these equations state: Intuitively this means that in 2/3 of the cases we initially pick a goat and the entertainer shows us where the car is by revealing the other goat. Only in 1/3 of the cases we picked the car initially and by changing our decision we pick a goat. 

Puzzles | Action sequences | Sharing Milk | Solution 

Puzzles | Action sequences | Crossing the River | Solution Label the farmer F, wolf W, goat G, cabbage C, river . and boat &lt; &gt;.  FWGC &lt; &gt;...  W C ..&lt;FG&gt;..  W C ...&lt; &gt; F G  W C ..&lt;F &gt;.. G  FW C &lt; &gt;... G  C ..&lt;FW&gt;.. G On this trip, the farmer can take either the Wolf or the cabbage. All else  ... remains the same.  C ...&lt; &gt; FWG  C ..&lt;FG&gt;.. W  F GC &lt; &gt;... W  G ..&lt;FC&gt;.. W  G ...&lt; &gt; FW C  G ..&lt;F &gt;.. W C  F G &lt; &gt;... W C  ..&lt;FG&gt;.. W C  ...&lt; &gt; FWGC 

Puzzles | Action sequences | Crossing the River II | Solution Label the fathers A, B and C, the sons a, b and c, river . and boat &lt; &gt;.  AaBbCc &lt; &gt;...  A B Cc ..&lt;ab&gt;..  A B Cc ...&lt; &gt; a b  A B Cc ..&lt;b &gt;.. a  A BbCc &lt; &gt;... a  A B C ..&lt;bc&gt;.. a  A B C ...&lt; &gt; a b c  A B C ..&lt;c &gt;.. a b  A B Cc &lt; &gt;... a b  Cc ..&lt;AB&gt;.. a b  Cc ...&lt; &gt; AaBb  Cc ..&lt;Bb&gt;.. Aa  BbCc &lt; &gt;... Aa  b c ..&lt;BC&gt;.. Aa  b c ...&lt; &gt; AaB C  b c ..&lt;a &gt;.. A B C  a b c &lt; &gt;... A B C  c ..&lt;ab&gt;.. A B C  c ...&lt; &gt; AaBbC  c ..&lt;C &gt;.. AaBb  Cc &lt; &gt;... AaBb  ...&lt;Cc&gt;. AaBb  ...&lt; &gt; AaBbCc 

The sequence is an = A(n,n), starting at a0. (A is of course the Ackermann function.) So the next number in the sequence is formula_1. 

Puzzles | Geometric puzzles | Matchsticks You have 6 unit length matchsticks, and you urgently need to make 4 triangles with unit length sides. You can't use anything except matchsticks to make the triangles, and the matchsticks are indestructable and non-replicable. How do you solve the dilemma? Solution 

Arrange the six matchsticks along the edges of a , or triangular pyramid. The matchsticks will form 4 triangular faces. 

Puzzles | Logic puzzles | Three Knights A man was sentenced to death, but the king wanted to give him a last chance. He asked the man to choose from one of the three knights that were there. One of the knights was the Knight of Life, and he always told the truth. The second Knight was the Knight of Death, and he always told lies. The third knight was the Knight of the Dungeon. He sometimes lied and sometimes told the truth. If the man chose the Knight of Death, he would be executed before sunset. If he chose the Knight of Life, he would be acquitted and set free right away. If he chose the Knight of the Dungeon, he would spend the rest of his life imprisoned in the Dungeon. The man was allowed to ask the three knights one question each. Thus, he asked the fat knight, "What is the name of this tall knight?" The reply was, "He is the Knight of Life." He asked the small knight, "What is the name of this tall knight?" The reply was, "He is the Knight of Death." Then he asked the tall knight, "Who are you?" "I am the Knight of the Dungeon" was the reply. After that, the man was able to correctly choose the Knight of Life, and was set free immediately. "Who is the Knight of Life, and who are the other two knights?" Solution 

Puzzles | Logic puzzles | Three Knights | Solution The small knight is the Knight of Life. The fat knight is the Knight of the Dungeon. The tall knight is the Knight of Death. Reasoning. The tall knight is not the Knight of the Dungeon, since if he were neither of the other knights could be the Knight of Life. Since he lied, he is not the Knight of Life. Therefore, the tall knight is the Knight of Death. The small knight must then be the Knight of Life, leaving the fat knight as the Knight of the Dungeon. Provided the fat knight tells the truth he can not be the Knight of Life as he is pointing out the tall knight as Knight of Life. The answer of the tall knight contradicts the answer of the fat knight. Therefore, the fat knight was lying. The tall knight can not be the Knight of Life, because he would have said so when directly asked, who he is. We can conclude that the small knight is the Knight of Life and he is pointing out the tall Knight as Knight of Death as he is always telling the truth. The last job position left for the fat knight is the Knight of Dungeon. The convict can be happy that the tall knight did not answer: 'Yes, I am the Knight of Life'. In this context, it would have been possible and would have left our convict very unfortunate, because his chance to find the Knight of Life would have been reduced by 50% A better tactic would be to ask each Knight: "If I asked you who the Knight of Life was, whom might you identify?" The Knight of Life would correctly identify himself. If asked "Who is the Knight of Life?" the Knight of Death would incorrectly identify either himself or the Knight of the Dungeon. Hence, when asked how he might respond, he would lie again, and state that he would have identified the correct Knight of Life (double negative). The Knight of the Dungeon is the unknown quantity, but when at least two Knights agree on someone, that is the correct identification. Unfortunately not foolproof, as if he would have identified the Knight of the Dungeon for his first lie, he could then pick to instead identify the Knight of Death for his second lie. 

__NOEDITSECTION__ Presentation.  Content and Contributions. This is, to the best of our knowledge, the world's first open content US History Textbook. The users are invited to tweak and refine this book until there is nothing better available. The authors are confident that this will happen because of the success of the Wikipedia site. 

This book describes Python, an open-source general-purpose interpreted programming language available for the most popular operating systems. The current versions are 3.x while versions 2.x are no longer supported, since 2020. This book describes primarily the versions 3.x, but does at times reference versions 2.x. There are a few implementations for Python 3 (and older): the standard implementation written in C, and PyPy, a JIT-compiled version written in RPython - a subset of Python. For Python 2 only there are Jython written in Java and IronPython written in C# for the .NET environment. 

The answer is 76. Rule:  T T=L²+LR+1 mod 100  L R 

The answer is 68. Rule:  T T=L²-LR+1 mod 100  L R 

Puzzles | Logic puzzles | Pizzas  Aaron, Betty, Charlie, Debbie, and Eric each ordered a pizza with three of the following five toppings: green pepper, mushrooms, onions, pepperoni, and sausage. The only topping that Aaron and Charlie had in common was sausage. The only topping Debbie and Eric had in common was pepperoni. The only topping Charlie and Betty had in common was mushrooms. The only topping Debbie and Betty had in common was green pepper. "Which toppings did each of them have on their pizza?" Solution 

Puzzles | Logic puzzles | Pizzas | Solution Aaron - green pepper, onions and sausage Betty - green pepper, mushrooms and onions Charlie - mushrooms, pepperoni and sausage Debbie - green pepper, pepperoni and sausage Eric - mushrooms, onions and pepperoni Reasoning. We are informed that B and D had green pepper, B and C had mushrooms, D and E had pepperoni and A and C had sausage, all marked with a Y.  | A | B | C | D | E  Green Pep | | Y | | Y |  Mushrooms | | Y | Y | |  Onions | | | | |  Pepperoni | | | | Y | Y  Sausage | Y | | Y | | We are also informed that the pairs B and D, B and C, D and E, and A and C, didn't share anything, except what is already marked with a Y. This eliminates the possibilities marked with an N.  | A | B | C | D | E  Green Pep | | Y | N | Y | N  Mushrooms | N | Y | Y | N |  Onions | | | | |  Pepperoni | | N | | Y | Y  Sausage | Y | N | Y | | Knowing everyone had 3 toppings allows us to place another Y.  | A | B | C | D | E  Green Pep | | Y | N | Y | N  Mushrooms | N | Y | Y | N |  Onions | | Y | | |  Pepperoni | | N | | Y | Y  Sausage | Y | N | Y | | And we know B didn't share with C or D.  | A | B | C | D | E  Green Pep | | Y | N | Y | N  Mushrooms | N | Y | Y | N |  Onions | | Y | N | N |  Pepperoni | | N | | Y | Y  Sausage | Y | N | Y | | Allowing us to place 2 more Ys.  | A | B | C | D | E  Green Pep | | Y | N | Y | N  Mushrooms | N | Y | Y | N |  Onions | | Y | N | N |  Pepperoni | | N | Y | Y | Y  Sausage | Y | N | Y | Y | D didn't share with E, A didn't share with C.  | A | B | C | D | E  Green Pep | | Y | N | Y | N  Mushrooms | N | Y | Y | N |  Onions | | Y | N | N |  Pepperoni | N | N | Y | Y | Y  Sausage | Y | N | Y | Y | N This allows us to place the remaining 4 Ys.  | A | B | C | D | E  Green Pep | Y | Y | N | Y | N  Mushrooms | N | Y | Y | N | Y  Onions | Y | Y | N | N | Y  Pepperoni | N | N | Y | Y | Y  Sausage | Y | N | Y | Y | N Surprisingly, no one chose green pepper, onions and pepperoni. 

Authors' To-Do Lists Organic chemistry &gt; Organic textbook expanded contents Foreword Foundational concepts of organic chemistry Introduction to reactions Overview of Functional Groups Alkenes Alkynes Haloalkanes Alcohols Chirality Dienes Aromaticity Aromatic reactions Spectroscopy Organometallics Chemistry of various functional groups Organic chemistry and biochemistry Periodic table Lewis structures Nomenclature Isomers Structure and properties Reactions Chirality 

Puzzles | Arithmetical puzzles | Luxury Cars A major luxury car company has three outlets on the east coast (in New York, Boston, and Miami) and two outlets on the west coast (in Los Angeles and San Francisco). The sales in Los Angeles were one quarter of the total sales plus one car. The sales in New York were one quarter of the rest (not counting the cars sold in L.A.) plus one car. The sales in San Francisco were one quarter of the rest (not counting L.A., N.Y.) plus one car. And the sales in Boston were one quarter of the rest (not counting L.A., N.Y., S.F.) plus one car. The balance so far showed that L.A. and San Francisco sold 100 cars more than N.Y. and Boston. However, with the remaining cars sold in Miami it is likely that the total sales of the east coast will be greater than the total sales of the west coast. "How many cars were sold in each city?" solution 

Puzzles | Geometric puzzles | Connecting Utilities Here is a problem; you are commissioned to connecting utilities to the three houses A, B, and C below. The utilities are G for gas, W for water, and E for electricity.  |A| |B| |C|  [G] [W] [E] You must connect each house to each utility without the connections crossing at all. Solution 

Puzzles | Geometric puzzles | Connecting Utilities | Solution This problem can be analyzed using graph theory. The problem is essentially showing that the bipartite graph K3,3 is planar. However, "Kuratowski's theorem" tells us that this graph must "not" be planar. Essentially, there is no solution and the required construction "cannot" be done! Sorry! :) However, if we look at where the three houses and utilities are on a "non-planar" surface, such as a "torus" (or doughnut), we obtain some topological niceties that allow us to solve this problem. Here is an example of one solution. Lines moving off the torus and looping around are signified with a V.  ==V===V=V==V====V=====V===V======  : | /-\' | | /-\' | /-\' :  : | |A| | | |B| | |C| :  : |__/ | | \_| |___/ | | :  : | \____________/ | :  : \______ ________/ :  : | | :  : [G]---\ [W] [E] :  ======V=V==V=====V=====V==V==V=== Here is a picture with forks, essentially equivalent to parallel connections The solution usually given to this puzzle depends upon the fact that the puzzle, as stated, does not prohibit one of the connections going under a house (for example, the gas connection for A is routed G-&gt;pass under B-&gt;A). This appears to be topologically equivalent to the above(?). 

^ Polish ^ As an initiation to Polish it's useful to know a few common expressions: ^ Polish ^ ] 

^ Polish ^ This question is used to ask about things/objects and refers mostly to both plural and singular, although the answers vary depending on the number: eg.  " - Co to jest?"  " - To jest książka." / "to yest kshõshka" / &gt; This is a book.  " - Co to jest?"  " - To są książki." / "to sõ kshõshki" / &gt; These are (some) books. To get information about a person, you need to use "kto?" / who? /: The possible answers are: eg.  " - Kto to jest?"  " - To jest człowiek." /"to yest chwovyek"/ This is a man/human.  " - Kto to jest?"  " - To są ludzie." /"to sõ looje"/ These are people. The following sentences may be employed to introduce your family or friends to someone:  ^ Polish ^ 

^ Polish ^ Polish pronunciation is rather regular. Once you learn the rules, you should be able to guess how a word is pronounced and get it more or less right even if you've never heard it before (unlike English which is rather unpredictable). Vowels are pronounced similarly to their counterparts in most other European languages (not English though) but note, there are no long vowels. Stress is almost always on the penultimate (next-to-last) syllable. Special letters are: Special letter combos are: Devoicing is not something you need to focus on but you should be aware of it. See also. ^ Polish ^ 

Depending on the classification chosen, there are either three or five genders in Polish: In plural "męskoosobowy" (masculine-personal) and "niemęskoosobowy" (non-masculine-personal) are used for masculine personal and the remaining ones respectively. Unlike German, and more like Italian, it is usually possible to find out gender by looking at the noun ending and the meaning of the noun. Now some examples:  ^ Polish ^ 

Intro to oxidation reactions. Oxidation reactions involve electron loss in the moeity of interest. When iron rusts, iron atoms lose electrons to oxygen molecules, and is said to undergo oxidation. Ozonolysis. This is one of the simplest and easiest reaction mechanisms to learn. In ozonolysis, a carbon-carbon double bond is split in two by reaction with ozone. The two products, which can be aldehydes or ketones, each have one oxygen atom double bonded to each of the two carbons that had made up the double bond. The ozone (O3) cuts right thru the carbon-carbon double bond and leaves oxygens double bonded to the two newly severed ends. that is, Osmylation. In this reaction, a carbon-carbon double bond turns into a single bond with one hydroxy group attached to each of the two carbons. When heated, this product splits into two ketones (each of the two carbons ends up double bonded to the oxygen in the hydroxy group). Epoxidation (peroxidation). This reaction uses peroxides (an oxide with an extra linked oxygen) to produce epoxides. 

Added links 

Combustion is the simplest and most basic organic reaction, so simple that it is often overlooked or taken for granted in organic chemistry texts and classes. Combustion is the breakdown of a compound in combination with oxygen to form carbon dioxide and water. Water and carbon dioxide are both very stable and low-energy, and the atoms that make up these molecules go to form them releasing energy. This energy is heat and light that is the source for campfires and internal combustion engines. 

The result of hydrogenation reactions is to increase the number of hydrogens on a molecule. The recipient molecule is typically an alkene whose double bond is reduced by the addition of a hydrogen to each of the two carbons. The mechanism of this reaction is usually through a catalyst, such as Nickel, Platinum (used in PtO2, otherwise known as Adams' Catalyst after Roger Adams) or Palladium (Pd) that is mixed onto an inert material like charcoal (sometimes referred to as Palladium on carbon). The reaction takes place not within a solution but actually specifically on the surface of the catalyst. First the H2 gas is absorbed onto the catalyst surface. The Alkene forms a complex with the catalyst surface due to the vacancy of the orbital on the metal interacting with the filled π orbital. Thus we have a structure that looks somewhat like this: Once the complex is formed the hydrogen is able to apply itself to the double bond twice and then the final product is fully saturated and floats away from the catalyst, thus regenerating it.Hydrocarbons are of two types, aliphatic and aromatic.There are different methods to identify a hydrocarbon, if no of C atoms:1 then Meth.eg. 

^Lesson 7^ Rellena los espacios en blanco - "Fill in the blank". 1. Colón __________ "(descubrir)" América en 1492. 2. Yo __________ "(dejar)" mi trabajo la semana pasada. 3. Anoche, los estudiantes __________ "(tener)" una gran fiesta. 4. Ayer, nosotros __________ "(ir)" a la playa. 5. El año pasado, Isabel y José se __________ "(casar)". 6. Ayer, Carlos no __________ "(venir)" a la universidad. 7. Anoche, yo __________ "(mirar)" la televisión. 8. Tú ________ "(asistir)" dos cursos esta mañana. 9. Alberto __________ "(comprar)" un carro el mes pasado. 10. ¡Felicitaciones! Vosotros __________ "(ganar)" el campeonato. Soluciones a los ejercicios ^Lesson7^ 

^ Indonesian ^ Kata Penghubung - "Conjunctions / Filler". Tapi - "But". Kata ini memiliki beberapa bentuk - "This has several forms" Kalau / Jika - "If / When". The words Kalau and Jika are interchangeable Meski / Meskipun - "Although". Synonymous to walau and walaupun. Jadi / Maka - "So/Therefore". The phrase maka dari itu is very often used to mean "therefore". Agar - "So that / In order to". Synonymous with "supaya". It is often used together "agar supaya" to mean "in order to", but using either one is fine. Saja - "Only / Just)". Synonymous with "hanya". However, if both are being used together as one phrase (i.e. in hanya saja), it means "however", or rather "the catch is". Lagi - "Again". If it is used in the front of the sentence, we should use the phrase sekali lagi or lagi-lagi to mean "again", or rather "one more time". Lagipula - "Moreover". &lt;hr&gt; ^ Indonesian ^ 

^ Indonesian ^ | Why Learn Indonesian? |How to use this Indonesian Wikibook » Why Learn Indonesian. Hi, welcome to this Indonesian tutorial. You might be wondering why on earth you should learn Indonesian. Allow me to persuade you. Demographic Reasons. Studying Indonesian means you can communicate with more than 240 million Indonesians, only a small percentage of whom are able to speak English. Bahasa Malaysia is also a close relative to Indonesian. You can understand both with ease since there are only minor differences in vocabulary. Therefore, learning Indonesian gives access to about 230 million people—including those in Malaysia, Singapore, and Brunei. Indonesian (or commonly called as 'Bahasa Indonesia' or simply 'Bahasa') has been taught in schools in Australia, the Netherlands, and Vietnam. In Timor-Leste, although Portuguese and Tetum are official languages, Indonesian is also important as working language as Timor-Leste was part of Indonesia from 1976 until their independence in 1999. Practical Reasons. Indonesian is derived from Malay, a language of South Sumatra which was broadly used for trade purposes in the Malay World (now Malaysia, Indonesia, Brunei, Singapore) for centuries. When Indonesians began their fight for independence from the Dutch, the Malay language was renamed "Bahasa Indonesia" (it was renamed in 1928, while Indonesia proclaimed its independence in 1945). In Malaysia the national language is referred to as "Bahasa Malaysia". Indonesian is thus a variant of Malay. Not even half of Indonesians are native speakers of the Indonesian language, especially those living in rural areas. Many have as a mother tongue one of a diverse array of local languages, including Javanese, Balinese, Sundanese, Madurese, Buginese, Batak, Minangkabau, or various Chinese dialects. But most Indonesians can speak Indonesian at least as a second language; it is taught in schools and understood even in the most remote islands. What a remarkably versatile language Indonesian is! It is the language of education across the country, from the primary school to the university. It is the language of government and business administration, media, literature, and of everyday life in the big cities. It is a must for foreigners living in Indonesia. And if you are on a business trip to Jakarta, or on vacation in Bali, knowing some Indonesian can really enrich the experience. For those who are just curious language learners or those with a scholarly bent, Indonesian has an immense collection of literature. Linguistic Reasons. Indonesian is very easy—honest! Learning it is a valuable experience in itself, and what's more: you can pick up the basics within a few weeks. Here's why it is easy: I hope now you can see why Indonesian is worth learning. The Catch. Now, the catch is that every language has a culture attached to it. Indonesian is no exception. Since the way Indonesian people think differs from most westerners, there are some hurdles in learning it. For example, most western people prefer active sentences, while Indonesians usually prefer passive sentences and omit the subject if it is not important. Also, in spoken Indonesian, the formal grammatical rules are often broken by lots of shortcuts, usually specific to the region, not to mention slang words and idioms. However, all Indonesians that have finished grade school should be able to speak and comprehend proper Indonesian. 

» Next section: Modular Arithmetic 

The bond dissociation energy, or bond enthalpy, for a diatomic molecule X-Y is defined as the energy required to break one mole of X-Y bonds, as illustrated in the following process... Bond enthalpies always refer to breaking bonds under gaseous conditions. The mean molar bond enthalpy is an average value that is quoted for a bond that can occur in different molecular environments. An example is methane, CH4 Bond enthalpy values are used in Hess's Law Calculations. The standard enthalpy of a reaction can be found by considering the bond enthalpies of the products and reactants of the reaction - For stronger bonds, bond dissociation energy is higher as more energy is needed to break the bond. A carbon-carbon double bond is stronger than a single bond and requires more energy to be broken. However, a carbon-carbon double bond is not twice as strong as a single one, it is only 1.5 times stronger. All chemical bonds need an input of energy to be broken, as bonds allow a lower energy state for the component atoms. If a bond did not offer a lower energy state for the atoms that form it, a bond would not form. 

Chemical equilibria are ratios relating the forward and backward direction of a reaction to each other. This ratio is represented by the letter K in the following equation: K = products / reactants Rate of reaction. Definitions. Rate of reaction is the speed at which a chemical reaction takes place, expressed as moles per unit time and unit volume. The rate "r" of a general reaction formula_1 is defined by: formula_2 from the expression above, it is clear that the usual convention is that reaction rate is taken as products formation rate. Rate is a function of the concentration of reactants and products, temperature, pressure and presence of a catalyst. Rate expression. A common expression for reaction rate is the "power law": formula_3 formula_4 is called the "kinetic constant", and formula_5, formula_6, etc. are called the "reaction order with respect to the reactant A, B" (or the "partial order of A, B" etc.), respectively. The sum of all the orders is the "global order" of the rate expression. Therefore the rate expression formula_7 is a "second-order" and "first-order" in A and B. The reaction orders are the same as the stoichiometric coefficient in the case of an elementary reaction only; in most cases they must be determined experimentally and are valid in the window of experimental conditions. Elementary data on reaction orders can be obtained by changing the concentration C of one of the reactant, say A, and measuring the initial rate r. Plotting a rate vs concentration log-log graph, one obtains a straight line whose slope is the partial order in—say—A by virtue of the relation formula_8 Influence of temperature and pressure. formula_9 as known as the equation Rate equations and reaction mechanisms. The most important research application of kinetic investigations is the determination of reaction mechanisms. In fact, the rate expression is "function" of it. From a postulated reaction mechanism (the "model"), a rate equation can be derived and used to analyse the experimental data. If the obtained fit is not statistically significant, the scheme is rejected. In complex systems, several schemes can produce compatible rate expressions and the problem of model discrimination is of primary importance. Limiting step. Often a reaction has two or more steps. One of the steps, usually the last one, is the slowest step, and is said to be "rate-limiting". Steady state approximation. Sometimes it is useful, when calculating the reaction rate, to assume that no particular step is rate-limiting. Instead, the reaction intermediate can either proceed to the product or return to the original reactant with an equal rate for either possibility. This is called the "steady-state approximation". Example. A classical example is the hydrolysis of haloalkanes: This reaction can occur by two mechanisms: the SN1 and SN. The former is a unimolecular substitution: its rate is determined only by the concentration of R-X, without regard to the concentration of the new substituent. The latter is bimolecular: its rate is first-order in both R-X and new substituent, for a combined rate order of 2. Equilibrium. Chemical equilibrium is the state when a net reaction is neither going forward nor backward. It is a "dynamic" equilibrium. The rate of the forward reaction equals the rate of the reverse reaction, so the two cancel each other out, and the "net" rate of change is zero. The chemical equilibrium is dictated by the equilibrium constant (often written "Keq"), expressed by the "mass-action law": formula_10 The concentration "C" can be expressed in any scale, e.g. molar franction, molarity, partial pressure. If the temperature and pressure are kept constant, no matter what are the initial concentrations, the system will evolve until the mass-action product is equal to the Keq. In general, systems with formula_11 are highly displaced to the reactant sides (almost no conversion at equilibrium), whereas when formula_12 the reaction goes to completion. From classical thermodynamics it can be showed that the following relation holds: formula_13 where formula_14 is the total change in Gibbs free energy with reaction (product minus reactants). Influence of temperature. The influence of temperature can be obtained by differentiation of the equation above to lead: formula_15 as known as the " equation". Therefore, for exothermic reaction (formula_16) an increase in temperature will decrease the formula_17 quantity, leading to a lower "Keq" , conversely for an endothermic reaction. 

Energy diagrams are used to show the favorability of a reaction. They show how energy gained or lost in the different stages of a reaction and show which stages are the slow and fast steps (slow steps have high potential energy). We can also compare the energies of one reaction to another in order to see which reaction will be favored. Transition states are high peaks in an energy diagram. If the end of the diagram is lower than the beginning, the product of the reaction is more stable and/or lower in energy than the starting materials, and the overall reaction is energetically favorable. If the tail end of the energy diagram is higher than the front, then the product is less stable or energetically favorable than the starting materials, and the overall reaction is energetically unfavorable. Any high peaks in the diagram indicate difficult points to pass and will slow down the reaction. 

&lt;br&gt;"A simple alkene undergoing hydrohalogenation" Alkenes are electron-rich. Just look at those electrons all clumped together in the carbon-carbon double bond. They act as (they give up electrons). Hydrohalogenation follows Markovnikov's rule. That is, the carbon of the double bond that starts out more substituted receives the halogen (becomes more substituated) while the carbon of the double bond that starts out less substituted ends up with the hydrogen. &lt;br&gt;"He who has, gets: This is another alkene illustrating hydrohalogenation. Notice how the more highly substituated carbon ends up with the X-" The H+ is attracted to the electron-rich double bond and adds first, on the carbon with more hydrogens, leaving the other carbon with a positive charge (a carbocation). Then the X- reacts with the carbocation. Read about Vladimir Vasilevich Markovnikov on "Wikipedia". 

Introduction. A "series" is the sum of a sequence of terms. An "infinite series" is the sum of an infinite number of terms (the actual sum of the series need not be infinite, as we will see below). An "arithmetic series" is the sum of a sequence of terms with a "common difference" (the difference between consecutive terms). For example: is an arithmetic series with common difference 3, since formula_2, formula_3, and so forth. A geometric series is the sum of terms with a common ratio. For example, an interesting series which appears in many practical problems in science, engineering, and mathematics is the geometric series formula_4 where the formula_5 indicates that the series continues indefinitely. A common way to study a particular series (following Cauchy) is to define a sequence consisting of the sum of the first formula_6 terms. For example, to study the geometric series we can consider the sequence which adds together the first n terms: Generally by studying the sequence of partial sums we can understand the behavior of the entire infinite series. Two of the most important questions about a series are: For example, it is fairly easy to see that for formula_8, the geometric series formula_9 will not converge to a finite number (i.e., it will diverge to infinity). To see this, note that each time we increase the number of terms in the series, formula_9 increases by formula_11, since formula_12 for all formula_8 (as we defined), formula_9 must increase by a number greater than one every term. When increasing the sum by more than one for every term, it will diverge. Perhaps a more surprising and interesting fact is that for formula_15, formula_9 will converge to a finite value. Specifically, it is possible to show that Indeed, consider the quantity Since formula_19 as formula_20 for formula_15, this shows that formula_22 as formula_20. The quantity formula_24 is non-zero and doesn't depend on formula_6 so we can divide by it and arrive at the formula we want. We'd like to be able to draw similar conclusions about any series. Unfortunately, there is no simple way to sum a series. The most we will be able to do in most cases is determine if it converges. The geometric and the telescoping series are the only types of series we can easily find the sum of. Convergence. It is obvious that for a series to converge, the formula_26 must tend to zero (because sum of an infinite number of terms all greater than any given positive number will be infinity), but even if the limit of the sequence is 0, this is not sufficient to say it converges. Consider the harmonic series, the sum of formula_27, and group terms This final sum contains "m" terms. As "m" tends to infinity, so does the sum, hence the series diverges. We can also deduce something about how quickly it diverges. Using the same grouping of terms, we can get an upper limit on the sum of the first so many terms, the "partial sums". or and the partial sums increase like formula_31, very slowly. Comparison test. The argument above, based on considering upper and lower bounds on terms, can be modified to provide a general-purpose test for convergence and divergence called the "comparison test" (or "direct comparison test"). It can be applied to any series with nonnegative terms: There are many such tests for convergence and divergence, the most important of which we will describe below. Absolute convergence. Theorem: If the series of absolute values, formula_38, converges, then so does the series formula_39 We say such a series "converges absolutely". Proof: Let formula_40 According to the Cauchy criterion for series convergence, exists formula_41 so that for all formula_42 : We know that: And then we get: Now we get: Which is exactly the Cauchy criterion for series convergence. formula_47 "The converse does not hold." The series formula_48 converges, even though the series of its absolute values diverges. A series like this that converges, but not absolutely, is said to "converge conditionally." If a series converges absolutely, we can add terms in any order we like. The limit will still be the same. If a series converges conditionally, rearranging the terms changes the limit. In fact, we can make the series converge to any limit we like by choosing a suitable rearrangement. E.g., in the series formula_48, we can add only positive terms until the partial sum exceeds 100, subtract 1/2, add only positive terms until the partial sum exceeds 100, subtract 1/4, and so on, getting a sequence with the same terms that converges to 100. This makes absolutely convergent series easier to work with. Thus, all but one of convergence tests in this chapter will be for series all of whose terms are positive, which must be absolutely convergent or divergent series. Other series will be studied by considering the corresponding series of absolute values. Ratio test. For a series with terms formula_26, if then E.g., suppose then so this series converges. Integral test. If formula_58 is a monotonically decreasing, always positive function, then the series converges if "and only if" the integral converges. E.g., consider formula_61, for a fixed formula_62. The integral converges, for formula_65, so the series converges. We can prove this test works by writing the integral as and comparing each of the integrals with rectangles, giving the inequalities Applying these to the sum then shows convergence. Limit comparison test. Given an infinite series formula_34 with positive terms only, if one can find another infinite series formula_32 with positive terms for which for a positive and finite formula_72 (i.e., the limit exists and is not zero), then the two series either both converge or both diverge. That is, "Example:" For large formula_6, the terms of this series are similar to, but smaller than, those of the harmonic series. We compare the limits. so this series diverges. Alternating series. Given an infinite series formula_34, if the signs of the formula_26 alternate, that is if for all "n" or for all formula_6, then we call it an "alternating series." The alternating series test states that such a series converges if and (that is, the magnitude of the terms is decreasing). Note that this test "cannot" lead to the conclusion that the series diverges; if one cannot conclude that the series converges, this test is inconclusive, although other tests may, of course, be used to give a conclusion. Estimating the sum of an alternating series. The absolute error that results in using a partial sum of an alternating series to estimate the final sum of the infinite series is smaller than the magnitude of the first omitted term. Geometric series. The geometric series can take either of the following forms As you have seen at the start, the sum of the geometric series is Telescoping series. Expanding (or "telescoping") this type of series is informative. If we expand this series, we get: Additive cancellation leaves: Thus, and all that remains is to evaluate the limit. There are other tests that can be used, but these tests are sufficient for all commonly encountered series. 

Appendices. Authors  __NOEDITSECTION__ 

This book, Introduction to Philosophy, was initiated by RMK (who also plans to complete it). The original version of the book was begun on October 22, 2003, making it the first wikibook about philosophy. Of course, this is just an outline at the moment, but "the journey of a thousand miles begins with a single step"... 

"Philosophy" is a word with numerous vastly differing definitions, ranging broadly and not always compatible with each other. Today, it is perhaps most often thought of as meaning an individual's set of guiding principles, mostly moral, that he refers to in planning out and living his life. However, philosophy as an intellectual or academic pursuit has little do with this, meaning rather something along the lines of: "The directed search for knowledge and systems of knowledge that explain topical phenomena such as the nature of existence, the causes of existence, the nature and causes of an individual, the nature and causes of knowledge itself, and a million other things." What these million other things are constitutes the essence, or "feel", of philosophy as a subject, and of what differentiates it from other kinds of directed research such as physics, biology, or even music -- all of which were at one point considered aspects of philosophy itself. As you read on, you'll develop a sense of these questions as they've evolved in the Western canon, and of how philosophy fits in with the entirety of human intellectual pursuits. Western philosophy as we know it today is generally considered to have originated in Ancient Greece, from the people's cosmogonical or "world-creation" myths. Thales of Miletus brought back kernels of knowledge from Ancient Egypt, which over time the Greeks developed into the earliest philosophy that remains recognizably such today. It is not surprising, then, that the term "philosophy" finds its roots in the Greek language. "Philo-" stems from the Greek word "philein", meaning 'to love', and "-sophy" comes from the Greek word "sophia", or wisdom. Philosophy, then, can be thought of as "the love of wisdom". Consequently, philosophers concern themselves with exploring some of the biggest questions that face humankind: What exists? What is real? What is the nature of the universe? Does God exist? Do humans have souls? What is knowledge? What is truth? How should individuals govern themselves? What is the nature of consciousness? Like many other fields of inquiry, philosophical questions spring from a certain kind of curiosity about life and reality. Much of philosophy confronts very abstract ideas and thus, unlike in other academic disciplines such as the sciences, philosophers cannot always rely on gathering empirical data through experimentation. Of course, there are numerous other techniques available to the philosopher, including formal logic, thought experiments, rational dissertation, colloquy, phenomenological analysis, and more. [The pursuit of philosophical wisdom can be understood as a quest for two different kinds of wisdom. Aristotle characterized these as theoretical wisdom and practical wisdom. The first sort seeks understanding primarily for its own sake, whereas the latter is concerned with understanding in order to "do". The distinction is not unlike those in other fields. The physicist attempts to understand matter and motion simply for the sake of extending the limits human knowledge, but the engineer attempts to understand matter and motion in order to "make" something useful. The philosopher who attempts to understand the metaphysical nature of the universe is concerned with a more theoretical area of philosophical inquiry, while his or her colleague who tries to understand the proper way for businesses to behave ethically, is pursuing a much more practical subject matter. Philosophy that concerns itself with these sorts of practical questions is usually called "Applied Philosophy".] Philosophy is perhaps unique in its lack of limitation regarding subject matter. Almost every conceivable idea has been or will be tackled by philosophy, and among those that aren't, most are areas of philosophy that matured sufficiently to break off from the discipline proper and form a new field of study: physics was once called "natural philosophy." In addition, philosophers typically ask questions that touch on almost every other field of academic inquiry. In fact, one major role of philosophy is to provide the underpinning for just about every other field of study, everything from the rules governing research in that field to its relationship to the other fields of human endeavor. As a result, most universities offer courses in the "philosophy of" other academic disciplines, of interest to both philosophers and students of the respective field. Despite these far-ranging interests, there are some unifying features of philosophical inquiry... 

Zaitsev's (sometimes spelled "Saytzeff") rule: In elimination reactions, the major reaction product is the most substituted alkene. The most substituted alkene is also the most stable. 

 &gt; Glossary 

A Textbook on Five Levels The question arose early in the development of this textbook as to precisely who would be the target audience. Although intended to be a "beginning" textbook on German, many felt that the early lessons were too difficult for younger students with very limited or no experience with German and, perhaps more importantly, limited skills in English grammar. For this reason a textbook on three levels was conceived. Beginning German (Level I) puts more emphasis on building vocabulary around subject matter interesting and useful to young students. Basic German (Level II) emphasises grammar, and assumes a greater knowledge of English grammar more typical of an older high school or a college student. If you are just beginning to learn German or attempting to teach yourself, you may wish to try both approaches and see which works better for you, since some people require a strong structural approach to learning a new language while others find this "structure" only impedes progress by adding another layer of complexity. Intermediate German (Level III), which requires even more knowledge of English, is for college students, preferably for sophomores or juniors. With even more complex lessons, grammar and vocabulary comes Advanced German (Level IV), which with the most complex and difficult parts of the German language, is for late college students (Seniors) and college graduates. The last level, which is a review level, but also has cultural facts and the history of the German language, is Reviewed German. (Level V). An existing, separate text, German/Grammar, may eventually be merged into the lesson modules or developed into useful appendices as a grammar reference. At present, however, German Grammar is an expanding, significant contribution to the textbook; it provides an important reference on German language grammar rules useful to the student working through any of the three levels. The German Language. German ("Deutsch") is a member of the western group of the Germanic languages. It is spoken primarily in Germany, Austria, the majority of Switzerland, Liechtenstein, Luxembourg, the Südtirol (South Tyrol) region of Italy, the Opole Voivodship of Poland, the eastern part of Belgium, parts of Romania, the Alsace (Elsass) region of France and parts of Denmark. Additionally, several former colonial possessions of these countries, such as Namibia in Africa, have sizable German-speaking populations. There are German-speaking minorities in several eastern European countries including Russia, and in the United States as well as countries in South America like Brazil, Argentina and Chile. Over 120 million people speak German as their native language. German is the third most popular foreign language taught worldwide, and the second most popular in Europe. Continue reading about the German language. German and English. If you are an English speaker unfamiliar with German, you may be surprised to learn that English and German are closely related languages and share many words that are very similar. Such words are called cognates. This is particularly true for everyday words in English that are Anglo-Saxon (that is, Germanic) in origin. Consider the following list of English words followed by their German counterparts: Some German words have the same origin as their English counterparts but the meaning has changed: Of course, even words whose spelling is no different in English and German may be pronounced quite differently. But in reading German, you will see the connections between these languages, even in many of the "small" words (the above examples are all nouns). For example: Note also the general similarity of sentence structure with English. The only real difference in the German is that the verb is moved forward in the sentence. However, there are many German sentences in which a verb form is the last word in the sentence. Unfortunately, while German is perhaps the easiest "foreign" language for an English speaker to learn, meanings of words that are spelled similarly are not always identical. These "false friends" can be confusing for the beginner. Further, German is a more structured language than English, with a more complex grammar, and it will become apparent as you learn German that you will also learn more about English language structure than you might ever recall from your high school English classes. For a quick listing of similarities and differences between English and German, read the Introduction to Level I. Vocabulary and Grammar. In learning to read or speak any language with which you have minimal acquaintance (that is, are not a native speaker of), the two aspects to be mastered are vocabulary and grammar. Acquiring vocabulary is a "simple" matter of memorization. For the language(s) we learn as children, this process is so transparent that we have trouble conceiving of the importance of having a large vocabulary. By the age of conscious recognition of our communicating with others through speech, we have already learned the meaning of thousands of words. Even words we have trouble defining, we readily understand their use in conversation. This process can be "reactivated," as it were, by immersion in a second language: a method of learning a new language by moving to a place where that language is spoken and having to get around and live without use of one's native tongue. The student of German language, if not residing in a German-speaking environment, must put forth substantial effort to learning words, including their meaning, their pronunciation and their usage in common sentences. Be sure to "learn"—commit to memory—all of the vocabulary words in each lesson as they are presented. Early lessons have simple sentences because it is assumed that the student's vocabulary is limited. But throughout the text, more complex discourses (often as photo captions) are included to introduce the student to regular German in use. It may be helpful to translate these using a German-English dictionary (access to one is a must; see Appendix 5 for on-line options). Other sources of German, such as newspapers, magazines, web sites, etc., can also be useful in building vocabulary and developing a sense of how German words are put together. The German Wikipedia provides an ever expanding source of German language articles that can be used for this purpose. Further, a German version of the Wikibooks project—a library of textbooks in German—is available at German Wikibooks. German grammar is more complex than, but sufficiently similar to, English that "reading" German is possible with minimal vocabulary in the sense that the student should generally recognize the parts of a sentence. With a good dictionary or an online translator, an English speaker can usually translate a German sentence close to correctly. However, to accurately speak and understand German, you must learn how each word functions in a sentence. There are eight basic grammatical functions: case, gender, number, tense, person, mood, voice, and comparison. How words "signal" these functions is an important aspect of learning a new language. English speakers should know all of these functions and the signals used in English, but it is often the situation that you know perfectly well how to speak English, without understanding much about word-functions and signals. For this reason, this textbook incorporates considerable detail on grammar, including both English and German grammar. The reference book "English" at "Wikibooks" may be consulted for additional help. When we say German is more complex than English, what we really mean is that the signals used in German are different from and more numerous than those used by English. Pronunciation. A guide to the pronunciation of German is provided. You should become familiar with this page early on, and refer to it often. Nothing can replace learning a language from a native speaker, but the text is liberally sprinkled with audio files providing the student with valuable input from hearing spoken German. Analyze the spoken words carefully. The pronunciation guide can only closely, not exactly, convey how German words should be pronounced. And of course, German (like English) has a number of dialects distinguished by differences in pronunciation. Help in the pronunciation of individual words can be found by accessing the sound files of either of the online dictionaries, links to which are given in the German websites appendix. Layout of Lessons. This textbook is intended as a beginning course in the German language for English speakers. Early lessons emphasize conversational subjects and gradually introduce German grammatical concepts and rules. In addition, sound files accompany appropriate parts of each lesson. Although the basic lessons ("Grundlegende Lektionen") are presented at about the (US) high school level, beginners (including those attempting to learn German outside of a course structure) are expected to work through several basic lessons up to an indicated point, when review is suggested along with additional study. The basic way lessons go to other lessons is very simple and direct: Layout within Lessons. The following subheadings or categories are offered within the lessons (Level II and above): The Student and the Lesson. Each level of the text is designed to constitute a course of study in the German language. For any level selected, each lesson should be read thoroughly and mastered before moving on. Substantial text in German is included and the student should read all of it, not once, but multiple times. At Levels II and III, complete translations into English are included only in selected places. Most of this text must be translated by the student using his or her acquired vocabulary and the vocabulary presented at the bottom of each lesson. As the German text is read (preferably out loud), the student must succeed in gaining an understanding of the meaning of each sentence, and of the role each word plays in establishing that meaning. To the beginner, there will seem to be many words in a German sentence that are out of place or even redundant or unnecessary. These add subtleties to the language that will make sense eventually. But it is important to experience these subtleties from the very beginning. Continue to the level one introduction. 

LEVEL III: MITTLERE STUFE. Section 1 ~ "Bonn, Germany" Section 2 ~ "Innsbruck, Austria" Section 3 ~ "Bavaria, Germany" Section 4 ~ "Ruhrgebiet, Germany" 



Water can pull a halide off of an alkane and leave a hydroxyl group in its place. Water can ... 

^Lesson 1^ Matching Sentences. Match the sentences/questions (1 to 5) to their corresponding responses (a to e). Soluciones a los ejercicios "Solutions to exercices" ^Lesson1^ 

^ Indonesian ^ | « Lesson 6: Particles | Lesson 7: Introducing Yourself | Lesson 8: My Family » Dialog Pertama "(1st Dialog)". Dialog "(Dialogue)". KERRY: Hai! Nama saya KERRY. Siapa namamu? AUDRIE: Hai KERRY! Nama saya AUDRIE. Salam kenal. KERRY: Salam kenal juga. AUDRIE: Apa kabar KERRY? KERRY: Baik-baik, Anda? AUDRIE: Baik-baik juga. KERRY: Anda tinggal di mana? AUDRIE: Di GADING NIAS. Anda? KERRY: Saya juga! AUDRIE: Permisi, saya harus pulang ke rumah. KERRY: Baik, baik, hati-hati AUDRIE! AUDRIE: Haha terima kasih KERRY. KERRY: Selamat tinggal! AUDRIE: Selamat tinggal! Terjemahannya "(translation)": KERRY: Hi! My name is KERRY. What's your name? AUDRIE: Hi, KERRY! My name is AUDRIE. Nice to meet you. KERRY: Nice to meet you too. AUDRIE: How are you KERRY? KERRY: Good, and you? AUDRIE: Also Good. KERRY: Where do you live? AUDRIE: In GADING NIAS, you? KERRY: Me too! AUDRIE: Sorry, I've got to go home. KERRY: "Okay", be careful AUDRIE! AUDRIE: Haha thank you KERRY. KERRY: Goodbye! AUDRIE: Goodbye! Catatan "(Note)". This dialogue illustrates typical informal introductions. Note here that the dialog use "nama saya" to mean to "my name" and "namamu" to mean to "your name". You can use the phrase "namaku" as well to mean "my name". To refresh our memory, note that the informal possessive pronouns are "-ku", "-mu", and "-nya", for first, second, and third singular person. For more review, you can click here. The noun "nama" is the root of "namaku" and "namamu"; and the suffixes "-ku" and "-mu" adding the possessive information. Certainly, you can substitute "namamu" with "nama Anda" for more formal situations. The phrase "salam kenal" roughly means "nice to meet you". Note that unlike English, normally Indonesians don't say anything after the introduction and then carry on with the conversation. So, this phrase is not often used in introductions. It's up to you. The word "juga" means "too". The word "pergi" and "pulang" both can be translated as "to go" in English. Only "pulang" is strongly associated with home. "pulang" means "to go home". "pergi" has always been used to point out where would you go to, anywhere else but your own house. Dialog Kedua "(2nd Dialogue)". Dialog "(Dialogue)". KERRY: Selamat pagi, Pak! Perkenalkan, nama saya KERRY. Pak WAKI: Oh! Selamat pagi, KERRY! Nama saya WAKI. Apa kabar? KERRY: Baik-baik. Terima kasih. Terjemahannya "(translation)": KERRY: Good morning, Sir! Let me introduce myself, my name is KERRY. Mr. WAKI: Oh! Good morning, KERRY! My name is WAKI. How are you? KERRY: Good. Thank you. Catatan "(Note)". This is a formal introduction, in casual situation. The word "kenal" means "to know someone". In this dialog, we use the inflected form "perkenalkan", which in this dialog context means "let me introduce myself". It is actually the command form of "memperkenalkan", which means "to introduce". Don't worry about how the words are composed. This time, you can just consider it as a single word. The phrase "apa kabar" means "how are you". As we've already read from lesson 1, it literally means "what news". The phrase "terima kasih" means "thank you". Actually, it literally means "receive love". Dialog Ketiga "(3rd Dialogue)". Dialog "(Dialogue)". KERRY: AUDRIE, perkenalkan, ini Pak WAKI. AUDRIE: Pak WAKI, nama saya AUDRIE. Pak WAKI: Halo, AUDRIE! Salam kenal. Terjemahannya "(translation)": KERRY: AUDRIE, let me introduce you, this is Mr. WAKI. AUDRIE: Mr. WAKI, my name is AUDRIE. Mr. WAKI: Hello, AUDRIE! Nice to meet you. Catatan "(Note)". This dialogue is to introduce someone to someone else. You should introduce the older person to the younger one as a rule of courtesy, as demonstrated in the dialogue. The dialogue assumes that Mr. Waki is older than Audrie. The younger person must then respond by addressing the older one also for courtesy. You can follow it by stating your name again, like the example above. Or, you can just say "Hai, Pak Waki!" instead. If both people are roughly of the same age, you can choose either one. Note the usage of the word "perkenalkan". It is appropriate for both introducing yourself and introducing someone else. ^ Indonesian ^ | « Lesson 6: Particles | Lesson 7: Introducing Yourself | Lesson 8: My Family » 

Parallelograms. A parallelogram is a geometric figure with two pairs of parallel sides. Parallelograms are a special type of quadrilateral. The opposite sides are equal in length and the opposite angles are also equal. The area is equal to the product of any side and the distance between that side and the line containing the opposite side. Properties of Parallelograms. The following properties are common to all parallelograms (parallelogram, rhombus, rectangle, square) 

General :Biology . 

&lt; They are ascending powers of 2: 1 2 4 8 16 32 64 ... Note : 2 raise power 0 (i.e 2^0) equals 1 

Add 1 2 3 4... OR x(x-1)/2+1 : 1 2 4 7 11 16 22 29 ... formula_1: 1 2 4 7 12 20 33 ... Add the three previous numbers: 1 2 4 7 13 24 44 ... The numbers with an odd number of ones in their binary representation: 1 2 4 7 8 11 13 ... or 1 10 100 111 1000 1011 1101 ... A formula for the formula_2th number is: formula_3, formula_4, formula_5. 

Deoxyribonucleic acid (DNA) and ribonucleic acid (RNA) are the information storage molecules and working templates for the construction of proteins. Every living cell and virus encodes its genetic information using either DNA or RNA. It is a true marvel of evolution that the vast amount of information needed to produce a human being can fit inside cells. Friedrich Miescher first isolated DNA and RNA from used surgical bandages in 1869. A series of experiments done by Oswald Avery, Colin MacLeod, and Maclyn McCarty in 1944 defined DNA as the genetic information. They injected virulent encapsulated bacteria into a mouse. As it a result, it died. On the contrary, nonvirulent nonencapsulated bacteria did not kill the mouse. Also, heating the virulent encapsulated bacteria and injecting did not kill the mouse either. The surprise lied in mixing the heat-killed virulent encapsulated bacteria with nonvirulent nonencapsulated bacteria did kill the mouse. Avery's group kept isolating factors until they discovered what killed the mouse, which was the DNA. In conclusion, the nonpathogenic bacteria was made pathogenic by DNA transfer from a pathogenic stain. A.D. Hershey and Martha Chase's work in 1952 reconfirmed that DNA is the carrier of genetic information. Hershey and Chase's group conducted an experiment involving two phages. One of the phages was radioactively labeled with heavy phosphorus to highlight the DNA and heavy sulfur to highlight the protein. The two phages inserted their "genetic material" into the bacterial cell. Afterwards, these cells were "blended" to remove the phage heads. The resulting radioactively-labeled heavy phosphorus was found in the cell. This experiment showed that DNA alone was sufficient to make new viruses, and therefore, defined as the genetic information. Rosalind Franklin, Chargaff, and other chemists provided information that led to the discovery of DNA as well. Franklin provided crystallographic data that implied a helical structure. Chargaff's rule said that the adenine content was equal to the thymine content and cytosine was equal to guanine. Also the aromaticity of the electron-rich purine and pyrimidine rings allowed for tautomer shifts between the keto and enol form. The keto tautomer (lactam) dominated at physiological pH. Garnering all this information, James Watson and Francis Crick discovered the structure of DNA (double-helix) in 1953. Structure of DNA and RNA. DNA is a polymer of nucleotides. Nucleotides consist of a pentose sugar (2-deoxyribose for DNA and ribose for RNA), a phosphate group (formula_1), and a nitrogen-containing base. There are five nitrogen bases - adenine, guanine, thymine, cytosine and uracil. Individual nucleotides are connected by phosphodiester bonds to form polynucleotides. DNA exists as a double helix of two (2) strands of polynucleotides. According to the principle of complementarity nucleotides A (adenine) bases are bound with a hydrogen bridge to the T (thymine) or U (uracil in case of RNA) on the other thread and similarly C (cytosine) bind themselves to G (guanine). This principle allows for DNA and RNA replication and for limited possibilities of repairing the genetic information when one of the threads gets damaged. Difference between DNA and RNA. DNA is the permanent genetic information storage medium, found in the nucleus of most cells of most living organisms. RNA is, in the case of eukaryotes, the medium that transfers the genetic information from the nucleus to the cytoplasm where proteins are synthesized. The most basic structural difference between a DNA molecule and an RNA molecule is that DNA lacks an hydroxyl group at its 2' carbon while RNA has an hydroxyl group at its 2' carbon. There are three major types of RNA: Each tRNA contains an "anticodon" which is a series of 3 bases that complements 3 bases on mRNA. Additionally, by 2005, other types of RNA, such as siRNA (small, interfering RNA) and miRNA have been characterized. DNA structure is typically a double stranded molecule with long chain of nucleotides, while RNA is a single stranded molecule " in most biological roles" and it has shorter chain of nucleotide. Deoxyribose sugar in DNA is less reactive because of the carbon- hydrogen bonds. DNA has small grooves to prevent the enzymes to attack DNA. on the other hand, RNA contain ribose sugar which is more reactive because of the carbon-hydroxyl bonds. RNA not stable in its alkaline conditions.RNA has larger grooves than that of DNA, which makes RNA easier to be attacked by enzymes. DNA has a unique features because of its helix geometry which is in B-form. DNA is protected completely by the body in which the body destroys any enzymes that cleave DNA. In contrast the helix geometry of RNA is in the A- form; RNA can be broken down and reused. RNA is more resistance to be damage by Ultra- violet rays, while DNA can be damaged if exposed to Ultra-violet rays. DNA can be found in nucleus, while RNA can be found in nucleus and cytoplasm. The bases of DNA are A,T,C,and G, DNA is a long polymer with phosphate and deoxyribose backbone. In contrast, RNA bases are A,U, C,and G, RNA backbone contain ribose and phosphate. The job of DNA on the cell is to store and transmit genetic information. RNA job is to transfer the genetic codes that are needed for the creation of proteins from the nucleus to the ribosome, this process prevents the DNA from leaving the nucleus, so DNA stays safe. Plasmid. Plasmids are small, circular DNA molecules that replicate separately from the much larger bacterial chromosome can be use in laboratory to manipulate genes. Plasmids can carry almost any gene and they can passed it on from one generation of bacteria to the next generation, therefore plasmids are key tools for “gene cloning” which is the production of multiple identical copies of a gene carrying piece of DNA =DNA Replication ;Semi-conservative replication= In the process of DNA replication, as cells divide, copies of DNA must be produced in order to transfer the genetic information. In DNA replication, the strands in the original DNA separate, and then each of the parent strands make copies by synthesizing complementary strands. Enzyme called helicase will start The process of replication by catalyzing the unwinding of protein of the double helix. This enzyme can break the H bonds between the complementary bases. These single strands now act as templates for the synthesis of new complementary strands. The energy for the reaction can be provided when a nucleoside triphosphate bonds to a sugar (at the end of a growing strand), and 2 phosphate group are cleaved. Then a T in the template forms H bond with an A in ATP, and a G on the template forms H bond with CTP. After the entire double helix of the parent DNA is copied, the new DNA molecule will form . The new DNA molecule is consist of one strand from the parent (original) DNA and one strand from the daughter strand( newly synthesized strand). This process is called "semi-conservative replication". Recombination. Recombination refers to ways in which DNA modification leads to distinct gene expressions and functions. Deletions, insertions and substitutions are the most useful changes implemented for the synthesis of new genes. Recombination techniques and recombinant DNA technologies has made possible the modification and creation of new genes for specific proposes. These techniques helps clone, amplified and introduce new DNA into suitable cells for replication. Restriction enzymes and DNA ligase play a vital role in the production of recombinant DNA. Vectors are useful for cloning. Plasmids or viral DNA (lambda Phage) used to introduced a DNA sequence into a cell is refer as vector. The semi-conservative hypothesis was confirmed by both Matthew Meselson and Franklin Stahl in 1958. The parent DNA was labeled with "heavy nitrogen," the radioisotope of fifteen. This radioactive nitrogen was made through growing E. coli with ammonium chloride (also with "heavy nitrogen") environment. Directly after the labeling of "heavy nitrogen" to the parent DNA strand, it was transferred to a medium with only regular nitrogen (isotope of fourteen). The proposed question was "What is the distribution between the two nitrogens in the DNA molecules after successive rounds of replication?" Using density-gradient equilibrium sedimentation, this question was answered. In this centrifugation, similar densities were grouped together in the tube. Afterwards, an absorbance was applied with ultra-violet light to determine the density bands of all the different kinds of DNA in the solution. After one generation, a unique single density band was shown. It was not exactly the band of heavy nitrogen nor the regular band of regular nitrogen, but rather halfway between the two. This proved that DNA was not conservative but semi-conservative due to this nitrogen hybrid. One generation after this one, there were two nitrogen-14 molecules and two hybrid nitrogen molecules. The probability and statistics of this second generation reassures the semi-conservative hypothesis. As DNA strands can combine and form molecules, it can also be melted, resulting in separation or denaturation. It is important to note that the heat from the melting process does not disrupt the DNA structure, but rather DNA helicase takes energy from the heat, separating the DNA strands. The denaturation process occurs at about ninety-four degrees Fahrenheit where the double strands are converted to single stands. Absorbance readings were taken to reconfirm the denaturation process. When the DNA molecule is in its double-stranded form, it absorbs less UV light. In its single-stranded form, it absorbs more UV light because the single strand is more exposed than its double-strand counterpart. This is known as the hypochromic effect. In a PCR, polymerase chain reaction, it utilizes this denaturation capability in its primary step. Afterwards, it is then annealed at fifty to sixty degrees Fahrenheit. In this step, this is where the primers bind to their corresponding flanking DNA base sequence. In the last step, DNA polymerization, it requires a few ingredients: heat-resistant DNA polymerase, magnesium ion, buffer, and corresponding nucleotides. The DNA polymerase requires heat resistance because throughout the entire process, it undergoes a constant temperatures changes. The magnesium stabilizes the phosphate backbone. The buffer prevents drastic changes in pH, because if the pH lies near the extremes on the scale, it may denature the DNA so drastically that it will not recombine. This final step of DNA polymerization happens at the "melting temperature." The melting temperature is the temperature at which half of the DNA is single-stranded and the other half is double-stranded. At this optimal temperature of seventy-two degrees Fahrenheit, the DNA is amplified with its respective primers, replicating at an exponential rate. =DNA manipulation: human's effort to desire DNA and "read" DNA= Restriction endonucleases. What's endonuclease?--Endonucleases are enzymes that can recognize specific base sequences in double-helical DNA and cleave, at specific places, both strands of that duplex. Many prokaryotes contain restriction enzymes. Biologically, they cleave foreign DNA molecules. Because the sites recognized by its own restriction enzymes are metylated, the cells' own DNA are not cleaved. One striking properties of endonucleases is that they recognize DNA cleavage sites by "two-fold symmetry". In another word, within the cleavage site domain (usually 4 to 8 base pair), there exists a center, around which the cleavage site rotate 180 degree, the whole cleavage site would be indistinguishable. The following graph shows an example of how endonuclease recognize cleavage site by symmetry..After cleavage from several different endonulcease, a DNA sample would be sliced into many different duplex fragments, which are then separated by electrophoresis. The restriction fragments are then denatured to form single stranded DNA, which are then transferred to a nitrocellulose sheet, where they are exposed to a 32P-labeled single-stranded DNA probe. The single-stranded DNA are then hybridized by the radio active complementary strand. These duplex are then examined by autoradiography, which reveals the sequence of the DNA. This technique is named Southern plot, after Edwin South, who invented the technique. The following image shows the steps for creating a southern plot. DNA sequencing. DNA subunits are dGMP, dAMP, dCMP and dTMP. In a DNA molecule, these subunits are connected in a way that their phosphate on 5 prime carbon are bonded to another subunit's hydroxy group on 3 prime carbon. If there is no hydroxy group on 3 prime carbon, then the DNA polymeration will be terminated. Utilizing this idea, subunit analogies were invented, ddGMP, ddAMP, ddCMP and ddTMP. These analogies have the same structure as their parent molecule except that they have hydrogen on 3 prime carbon instead of hydroxyl group. In practice, target DNA template is put into four polymeration environments, dNTP+ddGMP or ddAMP or ddCMP or ddTMP, respectively.. Statically, fragments with different length terminated at base, G, A, C, T will be formed, respectively. Then these fragments will be separated by electrophoresis, which will tell the sequence of the target DNA. DNA synthesis. Although nature synthesize DNA from 5 prime carbon to 3 prime carbon, human synthesize DNA from 3 prime carbon to 5 prime carbon. The following monomer is commercially available for each of the base and is used as basic building block for DNA synthesis. It starts with the subunit with the first base bonded to resin by the 3 prime hydroxyl group. The second monomer with its own 5 hydroxyl protected by DMT (protecting group) was activated on its 3 prime phosphate, which is attached to the 5 prime hydroxyl group by nucleophilic substitution. Then the phosphite is oxidized by iodine to yield phosphate. Then DMT is removed to make the 5 prime hydroxyl available for the incoming of the third monomer. The graphical representation is shown below. Polymerase Chain Reaction (PCR). It is used for amplification of DNA. It includes three stages: DNA denaturation, oligonucleotide annealing, and DNA polymerization. In each cycle, the DNA duplex is denatured at first, then the primers for both strands (forward and reverse) bind to the desired sites and at last both strands polymerize to form the desired sequence. What needed in PCR are the DNA containing the desired sequence, primers for both strands, heat resistant DNA polymerase, dNTP and constantly changing temperature. The following image shows the detailed steps in PCR. =RNA Editing= Introduction. Adenosine deaminases (ADAR) are mRNA editing enzymes that alternate a double stranded RNA sequence by tuning it. When tuning the mRNA, ADARs retype the nucleotides of the mRNA at certain, specific points they wish to change. Because ADARs can change the sequence of mRNA, they can create or remove a splice site, delete or alter the meaning of a codon, and finally change the sequence of the RNA altogether. ADAR plays a vital part in editing mRNA and modulating the mRNA translational activity and is mainly found in the nervous system and is important in regulating the nervous systems. Lacking the ADAR gene is very detrimental to health. Function. &lt;br&gt; RNA splicing and RNA editing are very similar to one another. RNA splicing is more of a cut, copy, and past process while RNA editing is an alteration of one or more nucleotides. There are mainly two types of ADAR that alternate the sequence of a single nucleotide. The first type of ADAR retypes cytidine nucleotides into uridine nucleotides while a second type of ADAR retypes adenosine into inosine. &lt;br&gt;&lt;br&gt; mRNA consisting of many inosine nucleotides in its sequence is known to be very abundant in the nucleus and only mRNAs that leave the nucleus can be translated. In this case, ADAR is modulating the translational activities of mRNAs. In modulating the translational activities of the specified mRNA, it is also modulating the abundance of the certain proteins that are expressed in the edited mRNA. Also, increasing inosine can stabilize the structure of mRNA. Hence, ADAR can function as an mRNA stabilizer too. &lt;br&gt;&lt;br&gt;ADARs target the coding sequence, introns, and 5’ and 3’ untranslated reigions (UTR). ADAR targets the coding sequence by changing the nucleotide which then changes the geometry of the protein. Change in the geometry of the protein may yield to higher protein-protein affinity interaction. This high protein-protein affinity can improve biological reactions by facilitating those reactions. Alternation of the non-coding sites such as introns, 5’, and 3’ UTR is still unknown. The only thing that is known is that alternating the sequence within these sites can create new splice sites. Importance of ADAR. &lt;br&gt;Adar is mainly found in the nervous system because it plays a very important role in regulating the nervous system. One of Adars’ primary functions is its regulation of the nervous system. Adar regulates the nervous system by editing the glutamate receptors’ mRNA. These glutamate receptors construct the glutamate-gated ion channels which are important in transferring electrical signals between neurons. When ADARs edit gluR mRNA, they change the calcium permeability of these glutamate channels. They can either lower the calcium permeability for certain glutamate channel or increase calcium permeability for other glutamate channels. Regulation of calcium permeability is important to proper neurotransmission between neurons. Disease connecting to ADAR. &lt;br&gt;Adenosine deaminase Deficiency (ADD) is caused by the lack of ADAR gene. Hence ADARs editing is important for survival especially for human. Without ADAR, the body can not break down the toxin deoxyadenosine. The accumulation of this toxin destroys the immune cells making its host vulnerable to infection from bacteria and viruses. Seizures can also be caused by ADAR. Unedited version of gluR mRNA strongly causes seizures. When a seizure happens, the body mechanism increases the level of ADAR activity to edit gluR mRNA to prevent future seizures. DNA Replication:Initiation. &lt;br&gt;There are multiple mechanisms that have been proposed for the melting of the DNA strand and the subsequent unwinding of the double helix during the initiation process of DNA replication. Two recently proposed models are E1 and LTag. These two models function differently to attain the same overall goal of melting the DNA double helix due to differences in their overall structure. E1 model: In the E1 hexamer, there are six ß-hairpins in the central channel of the protein. They are arranged in a staircase-like structure. Two adjacent E1 trimers assemble at the origin point to melt the double-stranded DNA. A ring shaped E1 hexamer is then formed around the melted single-stranded DNA and the DNA is then pumped through the ring hexamer to from a fork that allows for the replication of the DNA. Ltag: The channel diameters of these hexamers vary between 13-17 angstroms. The hexamer responsible for melting has an assortment of planar ß-hairpins in the narrow channel. Two of these hexamers surround the origin point of the double-stranded DNA and squeeze the area together. This causes the melting of the dsDNA by forcing the breakage of base pairing. The two single-stranded DNA strands are then pumped into a larger channeland finally out through two separate Zn-domains. This allows for the replication of the DNA. The two mechanisms here are built solely on the structural knowledge of the proteins involved in the process and as such more experimental support is needed to confirm. There are not an excessive amount of helicases. Currently, work is being done in discerning the more simple helicases, and trying to find the different mechanisms that they may have. While there are some similarities, for example helicases all seem to have a hexameric structure, there are also differences as well in their structures that contribute greatly to the differences in their mechanisms. -This article refers to a recently published article in Curr Opin Struct Biol. published 2010, Sep. 24.: "Origin DNA melting and unwinding in DNA replication" by Gai "et al." 

Introduction. Logic is the study of the way we reason. In this chapter, we focus on the "methods" of logical reasoning, i.e. digital logic, predicate calculus, application to proofs and the (insanely) fun logical puzzles. Boolean algebra. In the black and white world of ideals, there is absolute truth. That is to say "everything" is either true or false. With this philosophical backdrop, we consider the following examples: &lt;br&gt; That is (without a doubt) true! That is also true. But what about: It is true! Although 1 + 1 = 3 is not true, the OR in the statement made the whole statement to be true if "at least one" of the statements is true. Now let's consider a more puzzling example The truth or falsity of the statements depends on the "order" in which you evaluate the statement. If you evaluate "2 + 2 = 4 OR 1 + 1 = 3" first, the statement is false, and otherwise true. As in ordinary algebra, it is necessary that we define some rules to govern the order of evaluation, so we don't have to deal with ambiguity. Before we decide which order to evaluate the statements in, we do what most mathematician love to do -- replace sentences with symbols.&lt;br&gt; Let "x" represent the truth or falsity of the statement 2 + 2 = 4.&lt;br&gt; Let "y" represent the truth or falsity of the statement 1 + 1 = 3.&lt;br&gt; Let "z" represent the truth or falsity of the statement 1 - 3 = -1.&lt;br&gt; Then the above example can be rewritten in a more compact way:&lt;br&gt; To go one step further, mathematicians also replace OR by + and AND by ×, the statement becomes: Now that the order of precedence is clear. We evaluate (y AND z) first and then OR it with x. The statement "x + yz" is true, or symbolically where the number 1 represents "true". There is a good reason why we choose the multiplicative sign for the AND operation. As we shall see later, we can draw some parallels between the AND operation and multiplication. The Boolean algebra we are about to investigate is named after the British mathematician George Boole. Boolean algebra is about two things -- "true" or "false" which are often represented by the numbers 1 and 0 respectively. Alternative, T and F are also used. Boolean algebra has operations (AND and OR) analogous to the ordinary algebra that we know and love. Basic Truth tables. We have all had to memorize the 9 by 9 multiplication table and now we know it all by heart. In Boolean algebra, the idea of a truth table is somewhat similar. Let's consider the AND operation which is analogous to the multiplication. We want to consider: where and "x" and "y" each represent a true or false statement (e.g. It is raining today). It is true if and only if both "x" and "y" are true, in table form: We shall use 1 instead of T and 0 instead of F from now on. Now you should be able to see why we say AND is analogous to multiplication, we shall replace the AND by ×, so x AND y becomes x×y (or just "xy"). From the AND truth table, we have: To the OR operation. "x" OR "y" is FALSE if and only if "both" "x" and "y" are false. In table form: We say OR is almost analogous to addition. We shall illustrate this by replacing OR with +: The NOT operation is not a "binary operation" like AND and OR, but a "unary operation", meaning it works with one argument. NOT "x" is true if "x" is false and "false" if "x" is true. In table form: In symbolic form, NOT x is denoted x' or ~x (or by a bar over the top of x). Alternative notations: and Compound truth tables. The three truth tables presented above are the most basic of truth tables and they serve as the building blocks for more complex ones. Suppose we want to construct a truth table for xy + z (i.e. x AND y OR z). Notice this table involves three variables (x, y and z), so we would expect it to be bigger than the previous ones. To construct a truth table, firstly we write down all the possible combinations of the three variables: There is a pattern to the way the combinations are written down. We always start with 000 and end with 111. As to the middle part, it is up to the reader to figure out. We then complete the table by hand computing what value each combination is going to produce using the expression xy + z. For example: We continue in this way until we fill up the whole table The procedure we follow to produce truth tables are now clear. Here are a few more examples of truth tables. Example 1 -- x + y + z  Example 2 -- (x + yz)' When an expression is hard to compute, we can first compute intermediate results and then the final result.  Example 3 -- (x + yz')w Exercise. Produce the truth tables for the following operations: Produce truth tables for: Laws of Boolean algebra. In ordinary algebra, two expressions may be equivalent to each other, e.g. xz + yz = (x + y)z. The same can be said of Boolean algebra. Let's construct truth tables for: xz + yz (x + y)z By comparing the two tables, you will have noticed that the outputs (i.e. the last column) of the two tables are the same! Definition We list a few expressions that are equivalent to each other "Take a few moments to think about why each of those laws might be true." The last law is not obvious but we can prove that it's true using the other laws: It has been said: "the only thing to remember in mathematics is that there is nothing to remember. Remember that!". You should not try to commit to memory the laws as they are stated, because some of them are so deadly obvious once you are familiar with the AND, OR and NOT operations. You should only try to remember those things that are most basic, once a high level of familiarity is developed, you will agree there really isn't anything to remember. Simplification. Once we have those laws, we will want to simplify Boolean expressions just like we do in ordinary algebra. We can all simplify the following example with ease: the same can be said about: From those two examples we can see that complex-looking expressions can be reduced very significantly. Of particular interest are expressions of the form of a "sum-of-product", for example: We can factorise and simplify the expression as follows It is only hard to go any further, although we can. We use the identity: "If the next step is unclear, try constructing truth tables as an aid to understanding." And this is as far as we can go using the algebraic approach (or any other approach). The algebraic approach to simplification relies on the principle of elimination. Consider, in ordinary algebra: We simplify by rearranging the expression as follows Although we only go through the process in our head, the idea is clear: we "bring" together terms that cancel themselves out and so the expression is simplified. De Morgan's theorems. So far we have only dealt with expressions in the form of a "sum of products" e.g. xyz + x'z + y'z'. De Morgan's theorems help us to deal with another type of Boolean expressions. We revisit the AND and OR truth tables: You would be correct to suspect that the two operations are connected somehow due to the similarities between the two tables. In fact, if you invert the AND operation, i.e. you perform the NOT operations on x AND y. The outputs of the two operations are almost the same: The connection between AND, OR and NOT is revealed by "reversing" the output of x + y by replacing it with x' + y'. Now the two outputs match and so we can equate them: this is one of de Morgan's laws. The other which can be derived using a similar process is: We can apply those two laws to simplify equations: Example 1&lt;br&gt; Express "x" in "sum of product" form Example 2&lt;br&gt; Express "x" in "sum of product" form Example 3&lt;br&gt; Express "x" in "sum of product" form Example 4&lt;br&gt; Express "x" in "sum of product" form Another thing of interest we learnt is that we can "reverse" the truth table of any expression by replacing each of its variables by their opposites, i.e. replace x by x' and y' by y etc. This result shouldn't have been a surprise at all, try a few examples yourself. De Morgan's laws Propositions. We have been dealing with propositions since the start of this chapter, although we are not told they are propositions. A proposition is simply a statement (or sentence) that is either TRUE or FALSE. Hence, we can use Boolean algebra to handle propositions. There are two special types of propositions -- tautology and contradiction. A tautology is a proposition that is always TRUE, e.g. "1 + 1 = 2". A contradiction is the opposite of a tautology, it is a proposition that is always FALSE, e.g. 1 + 1 = 3. As usual, we use 1 to represent TRUE and 0 to represent FALSE. Please note that opinions are not propositions, e.g. "42 is an awesome number" is just an opinion, its truth or falsity is not universal, meaning some think it's true, some do not. Examples. Since each proposition can only take two values (TRUE or FALSE), we can represent each by a "variable" and decide whether compound propositions are true by using Boolean algebra, just like we have been doing. For example "It is always hot in Antarctica OR 1 + 1 = 2" will be evaluated as true. Implications. Propositions of the type if "something" "something" then "something" "something" are called implications. The logic of implications are widely applicable in mathematics, computer science and general everyday common sense reasoning! Let's start with a simple example is an example of implication, it simply says that 2 - 1 = 1 is a consequence of 1 + 1 = 2. It's like a cause and effect relationship. Consider this example: There are four situations: In which of the four situations did John NOT fulfill his promise? Clearly, if and only if the second situation occurred. So, we say the proposition is FALSE if and only if John becomes a millionaire and does not donate. If John did not become a millionaire then he can't break his promise, because his promise is now claiming nothing, therefore it must be evaluated TRUE. If "x" and "y" are two propositions, "x" implies "y" (if "x" then "y"), or symbolically has the following truth table: For emphasis, formula_19 is FALSE if and only if "x" is true and "y" false. If "x" is FALSE, it does not matter what value "y" takes, the proposition is automatically TRUE. On a side note, the two propositions "x" and "y" need not have anything to do with each other, e.g. "1 + 1 = 2 implies Australia is in the southern hemisphere" evaluates to TRUE! If then we express it symbolically as It is a two way implication which translates to "x" is TRUE if and only if "y" is true. The "if and only if" operation has the following truth table: The two new operations we have introduced are not really new, they are just combinations of AND, OR and NOT. For example: "Check it with a truth table". Because we can express the "implication" operations in terms of AND, OR and NOT, we have open them to manipulation by Boolean algebra and de Morgan's laws. Example 1&lt;br&gt; Is the following proposition a tautology (a proposition that's always true) Solution 1&lt;br&gt; Therefore it's a tautology. Solution 2&lt;br&gt; A somewhat easier solution is to draw up a truth table of the proposition, and note that the output column are all 1s. Therefore the proposition is a tautology, because the output is 1 regardless of the "inputs" (i.e. x, y and z). Example 2&lt;br&gt; Show that the proposition z is a contradiction (a proposition that is always false): Solution Therefore it's a contradiction. Back to Example 1, :formula_24. This isn't just a slab of symbols, you should be able translate it into everyday language and understand intuitively why it's true. Logic Puzzles. Puzzle is an all-encompassing word, it refers to anything trivial that requires solving. Here is a collection of logic puzzles that we can solve using Boolean algebra. Example 1 We have two type of people -- knights or knaves. A knight always tell the truth but the knaves always lie. Two people, Alex and Barbara, are chatting. Alex says :"We are both knaves" Who is who? We can probably work out that "Alex" is a knave in our heads, but the algebraic approach to determine "Alex" 's identity is as follows: we simplify: Therefore "A" is FALSE and "B" is TRUE. Therefore Alex is a knave and Barbara a knight. Example 2 There are three businessmen, conveniently named Archie, Billy and Charley, who order martinis together every weekend according to the following rules: Putting all these into one formula and simplifying: formula_43 Exercises. Please go to Puzzles/Logic puzzles. Problem Set. 1. Decide whether the following propositions are equivalent: 2. Express in simplest sum-of-product form the following proposition: 3. Translate the following sentences into symbolic form and decide if it's true: 4. NAND is a binary operation: Find a proposition that consists of only NAND operators, equivalent to: 5. Do the same with NOR operators. Recall that x NOR y = (x + y)' 

Puzzles | Logic puzzles | 10 Hats Ten students are caught cheating on a math exam and the teacher wants to give them a final chance. He makes them sit on separate rows, each one behind the other, facing the board. Each one sees all others at the front (i.e. 10th row student sees 9 students, 9th sees 8 students, ..., 1st one sees no one). He puts a hat on each student's head randomly, white or red. If at least 9 of them guess what color hat is on his or her head, the teacher will release them without punishment. Each student is allowed to talk once and may only say one word, which is his or her guess. They may decide on a strategy before starting to talk. Once they begin, they are not allowed to say or do anything else but guessing a color. Only one student is allowed to make a mistake. "What kind of a strategy should they follow?" See also Infinite Hats. Solution 

Puzzles | Logic puzzles | 10 Hats | Solution = Solution 1 = The 10th student says "white" if she sees an even number of white hats and "red" if she sees an odd number of white hats (she will be the only one that could make a mistake). The 9th student now sees 8 hats in front of her. Say X of them are white. Now there are four possibilities: The 9th student can figure out which case occurs and tell her correct hat color. From then on the students have to subtract the correctly guessed hats from the total and apply the even/odd reasoning again (using what the 10th student has said and counting the number of white hats in front of them). = Solution 2 = There is a slight variation on this solution by transformation: instead of wearing red and white hats, the students wear hats labeled 0 or 1. Let x1...,x10 be the numbers on the students' hats. Let a1...,a10 be the students' answers. The students answer in decreasing order as in the first solution. The answers are defined recursively as These equations can be succinctly described as follows: each student sums up the numbers they see and the numbers they've heard. The student says one if the grand total is odd and zero otherwise. To prove that ai === xi mod 2, let si = x1+...+xi and use induction. Base case: a9 === x9 mod 2 "Proof:" Inductive step: Assume ai === xi mod 2 for i = n+1..., 9 Then Thus ai === xi mod 2 for i = n, n+1..., 9 Therefore ai === xi mod 2 for i = 1...,9 and at most one answer (a10) is incorrect. = Solution 3 = The 10th student will see 9 hats with only one or two colours. One of those colours has to have an odd number showing. The initial setup might be easier for the students if the 10th student just said which colour has an odd number showing. The rest of the students will be correct through the same logical reasoning as in Solution 1. 

The preposition means "at" or "in": The contraction is used in place of "à le" (singular): Likewise, the contraction is used in place of "à les" (plural). Mireille: Monique: Mireille: Monique: Marcelle: Similar to English, "pleuvoir" is an : it has only a third-person singular conjugation: In order to say that one did "not" do something, the construction must be used. The is placed before the verb, while the is placed after. Formation and rules. Simple negation is done by wrapping around the verb: In a past tense, surrounds the auxiliary verb, not the participle: When an infinitive and conjugated verb are together, usually surrounds the conjugated verb:  can also precede the infinitive for a different meaning:  precedes any pronoun relating to the verb it affects: In spoken French, the can be omitted, leaving simply after the verb in context: Negation of indefinite articles. The indefinite articles "un", "une", and "des" change to "de" (or "d’") when negating a sentence. Note that means both "the weather" and "the time". The verb is translated to "to go". It is irregularly conjugated (it does not count as a regular verb). Usage. There is no present progressive tense in French, so "aller" in the present indicative is used to express both "I go" and "I am going": "Aller" must be used with a place and cannot stand alone. In addition to meaning "at" or "in", the preposition means "to" when used with : An infinitive preceded by is used to say that something is going to happen in the near future: Recall that the negative goes around the conjugated verb. In place of a preposition and place, the pronoun , meaning "there", can be used; "y" comes before the verb: Remember that "aller" must be used with a place ("there" or a name) when indicating that you are going somewhere, even if a place wouldn't normally be given in English. The negative form of "aller" with the "y" pronoun has both the verb and pronoun enclosed between "ne" and "pas":  

 __NOEDITSECTION__ This book covers the C++ programming language, its interactions with software design and real life use of the language. It is presented in a series of chapters as an introductory prior to advance courses but can also be used as a reference book. This is an open work; if you find any problems with terms or concepts "you can help by contributing to it"; "your participation is needed and welcomed!" You are also welcomed to state any preference, shortcomings, vision for the actual book content, structure or other conceptual matters; see . Appendix A: References Tables. • Keywords • Preprocessors Directives • Standard Headers • Data Types • Operators • Standard C Library Functions • ASCII chart Appendix B: External References. • Weblinks • Books 

=Formulas= Computing factors of polynomials requires knowledge of different formulas and some experience to find out which formula to be applied. Below, we give some important formulas: formula_1 formula_2 formula_3 formula_4 formula_5 formula_6 Examples.  :formula_7  :formula_8 formula_9&lt;br&gt; formula_10&lt;br&gt; formula_11&lt;br&gt; formula_12  :formula_13  :formula_14  :formula_15  :formula_16 formula_17&lt;br&gt; formula_18&lt;br&gt; formula_19&lt;br&gt; formula_20&lt;br&gt; formula_21  :formula_22  :formula_23  :formula_24  :formula_25  :formula_26  :formula_27  :formula_28  :formula_29 formula_30&lt;br&gt; formula_31&lt;br&gt; formula_32&lt;br&gt; formula_33  :formula_34  :formula_35  :formula_36  :formula_37 write out the coefficients and if the end is equal to zero, than it is a root example: formula_38 formula_39 = Possible Factors = To factor we must first look for possible factors. Possible factors are any number that might be a factor. Once we have a possible factor then we divide that number into the number we are factoring. If they divide evenly then we have a factor! The factor is the possible factor we found and the result of the division problem. Here is an example. Let's say the number we are factoring is 20. 2 is the possible factor. 20 / 2 = 10. They divide evenly which means we have a factor. The factors are 2 (the possible factor), and 10 (the result of the division problem). Now that we have a factor we start over with a new possible factor and find all of the factors. Examples. Factor 12 First find all the possible factors The possible factors are 1, 2, 3, 4, 5, 6, 7, 8, 9, 10, 11, and 12 Next we will try them one by one 12/1 = 12 (1 and 12 are factors) 12/2 = 6 (2 and 6 are factors) 12/3 = 4 (3 and 4 are factors) 12/4 = 3 (we already have the factors 3 and 4) Once we get a factor we already have then we know we have all the factors. So the factors for 12 are 1, 2, 3, 4, 6, and 12. Factor 54 First find all the possible factors Do not worry this is not as much work as it seems! 54/1 = 54 (1 and 54 are factors) 54/2 = 27 (2 and 27 are factors) 54/3 = 18 (3 and 18 are factors) 54/4 = 13r2 (4 is not a factor) 54/5 = 10r4 (5 is not a factor) 54/6 = 9 (6 and 9 are factors) 54/7 = 7r5 (7 is not a factor) 54/8 = 6r6 (8 is not a factor) 54/9 = 6 (we already have the factors 9 and 6) So the factors for 54 are 1, 2, 3, 6, 9, 18, 27, and 54 Factor 180 First find all the possible factors Do not worry this is not as much work as it seems! 180/1 = 180 (1 and 180 are factors) 180/2 = 90 (2 and 90 are factors) 180/3 = 60 (3 and 60 are factors) 180/4 = 45 (4 and 45 are factors) 180/5 = 36 (5 and 36 are factors) 180/6 = 30 (6 and 30 are factors) 180/7 = 25r5 (7 is not a factor) 180/8 = 22r4 (8 is not a factor) 180/9 = 20 (9 and 20 are factors) 180/10 = 18 (10 and 18 are factors) 180/11 = 16r4 (11 is not a factor) 180/12 = 15 (12 and 15 are factors) 180/13 = 13r11 (13 is not a factor) 180/14 = 12r12 (14 is not a factor) 180/15 = 12 (we already have the factors 15 and 12) So the factors for 180 are 1, 2, 3, 4, 5, 6, 9, 10, 12, 15, 18, 20, 30, 36, 45, 60, 90, and 180. = Dividing polynomials = The process of factoring will require dividing polynomials. This form of division is not too different from long division method, and is known as "synthetic division". Consider the polynomial x3 - 21x2 + 143x - 315. In this case, determining factors may require trial and error (until you learn of alternate techniques), and when you do, you will need to divide the polynomial with the discovered factor. In this example, we will divide by (x-5). The full division starts like this:  1x -5 | 1x^3 -21x^2 +143x - 315 As with long division, you need to find the number used for subtraction and place it on top - in this case, you need to make sure the left-most term becomes zero. Next, multiply the newly added top-most term with the left hand side to get the amount to subtract, and perform the subtraction.  1x^2  1x -5 | 1x^3 -21x^2 +143x - 315  1x^3 -5x^2  -16x^2 +143x - 315 Repeat until the division is complete:  1x^2 -16x + 63  1x -5 | 1x^3 -21x^2 +143x - 315  1x^3 -5x^2  -16x^2 +143x - 315  -16x^2 + 80x  63x - 315  63x - 315  0 Some people may find writing the x3 and other variables to be bulky - if writing on pen and paper, they can be omitted as part of shorthand.  1 -16 + 63  1 -5 | 1 -21 +143 - 315  1 -5  -16 +143 - 315  -16 + 80  63 - 315  63 - 315  0 In this case, factoring is straight forward since you can easily determine the number to use for the next step in division. =See also= 

In a way, this book is about itself. One of the dominant trends in Constructivist theory stresses the social and cultural basis of human knowledge. is a prime example of this idea. So, Access to technology is in the schools, let's use constructivism with kids, even little kids to get them thinking and interested in learning for learning's sake. It needs to hit the grass roots and move once and for all out of the theoretical stages. Bridging, scaffolding, whatever you want to call it, children have to be exposed to this therefore educators have to understand the practical application. How about kids using blogs as well as wikis. Both seem to be naturals for constructivism. It is just up to teachers to figure out how to move from just posting a daily journal or daily assignments to considering the audience and engaging in responding, rethinking, analyzing. It has started at the university level and at the high school level; it is even happening in the classrooms of the younger kids. One thing though, 'defining proper use of blogs and wikis' , seems as if it would serve no purpose. It is important to know how blogs and wikis can be used, but it may not be constructive to define 'how'. Teachers, you are already using this theory, building from the known. Discussion and community is essential. Guided Reading and Making Words may be considered constructivist. Below is a section of this page for specific uses and development of constructivist theories. Classroom teachers will need practical applications. Foundations. Guided Reading. Guided Reading is by its nature a constructivist application. One moves from where the children are, what their existing knowledge is, into the new learning in the story. And we are already doing it. Making Words. Making Words starts at the knowledge of letters and sounds and evolves to an increasingly advanced knowledge of words. 

The Alphabet. Like English, the German alphabet consists of 26 basic letters. However, there are also combined letters and three umlauted forms. An "umlaut" is the pair of dots placed over certain vowels; in German, "Umlaut" describes the dotted letter, not just the dots. As in English, letters may be pronounced differently depending on word and location. The first column is the German letter, the second describes the IPA pronunciation and rough English approximation of the letter name. The third gives an English word that matches or approximates the German letter sound. Reading down this column and pronouncing the "English" words will recite the alphabet "auf Deutsch" ("in German"). Note that letter order is exactly the same as in English, but pronunciation is not for many of the letters. In the list of pronunciation notes, no entry means essentially "pronounced as in English". Deutsche Aussprache ~ German Pronunciation Guide. Vokale ~ Vowels. German vowels are either long or short, but never drawled as in some English dialects. A simple method of recognizing whether a vowel is likely to be long or short in a German word is called the Rule of double consonants. If a vowel is followed by a single consonant — as in "haben" (have), "dir" (you, "dat."), "Peter" (Peter), and "schon" (already) — the vowel sound is usually long. If there are two or more consonants following the vowel — as in "falsch" (false), elf" (eleven), immer" (always), and "noch" (still) — the vowel sound is usually short. There are some German words that are exceptions to the double consonant rule: "bin", "bis", "das", "es", "hat", and "was" all have short vowel sounds. It is also the case that the silent 'h' does not count as a consonant and the preceding vowel is always long, e.g. "ihnen". This "rule" is applied to the use of 'ss' vs. 'ß' (see below)—'ß' is treated as a single consonant for purposes of vowel length. Thus, the vowel before 'ß' in "der Fuß" /fuːs/ (foot) is long, while that before 'ss' in "das Fass" /fas/ (cask) is short. Vowel combinations. Note that 'ie' and 'ei' are pronounced in the opposite manner as they would be pronounced in English, where the rule is that the first vowel is long and the second is silent. In German, "die" is pronounced 'dee', but in English it sounds like 'dye'. The word "mein" in German is pronounced like the English 'mine'. A useful tip for English speakers learning German is to pronounce the English name of the second vowel in the combination. Konsonanten ~ Consonants. Most German consonants are pronounced similarly to the way they are pronounced in English, with exceptions noted in column 3 above. Details of certain consonant sounds and uses are discussed further here: German Sounds not found in English. There are sounds in the German language that have no real equivalent in the English language. These are discussed here. Audio: (37KB) ~ ach, auch, ich, richtig Syllable Stress. The general rule in German is that words are stressed on the first syllable. However, there are exceptions. Almost all exceptions are of Latin, French, or Greek origin. These words are generally stressed on the last syllable, for example, Vokal, Konsonant, and Lektion. These words (not stressed on the first syllable) appear in the (Level II and III) lesson vocabularies as "Vokal", "Lektion" (in some regions: "Lektion"), etc. Words starting in common prefixes (ge-, be-, ver-, etc.) stress the syllable following said prefix. Examples are Gemüse, Beamte, and Vereinigung. Links. For very advanced Readers: 

Tenses. As we learned from Lesson 4, Indonesian has no tenses. In order to express idea in different time frame, we need to attach time signals, such as "yesterday", "tomorrow", "this morning", etc. These time signals are very easy to learn. When you say a sentence without any time signal, we can never be sure what time frame it is assumed to be, if it is taken out of context. For example: Saya makan apel. The general translation would be in present tense: "I eat apple". However, it depends on the speaker on what it means. It may also mean progressive tense. The speaker may be eating an apple as he/she speaks. Also, you can incorporate as many time signals as you want to express more specific ideas, as long as the addition doesn't contradict the existing ones and follows the "general rule of thumb". This is especially useful since Indonesian has no notion of complex grammar such as future perfect. Present Tense. To express habitual activity, we use present tense. In English, we use the form of "infinitive + (-s/-es)". In Indonesian, we use the time signal "setiap X", where X can be subtituted with "hari", "pagi", "siang", "malam", "minggu", etc to denote that the activity is recurrent. Remember that when you mention a noun, it is uncertain whether it is singular or plural, as discussed here. The assumption is always singular, but it can still mean plurals. Here's a helpful word signals appropriate for present tenses: Note: The word setiap may be shortened as tiap. Both are acceptable in formal written/spoken Indonesian. To indicate that the action is a habit, you can put the word biasa right before the verb: Progressive Tense. To express a currently ongoing activity, we use progressive tense. In English, we use the form of to be + infinitive + -ing. In Indonesian, we use the time signal sedang. The word other than sedang that can be used is lagi. Be careful with the word placement. The word "sedang" or "lagi" are used right before the verb in order to form progressive tense. Be extra careful in the word "lagi", because if the placement is wrong, then it will mean "again" instead of to mean a progressive tense. You can also attach time signals to further reinforce your idea; such as: "sekarang" = now. Past Tense. Indonesian only has one notion of past tense, which is simple past. It has no notion of past progressive or past perfect tenses. (See below for further clarification) As always, to form the tenses, we just need to attach time signals. To express undefinite past, Indonesian has these phrases: Both the words sudah and telah literally means already. It explains that the action has already happened. It is uncertain whether it's in the recent or distant past. Note that due to English influence, sudah or telah are often used to express past perfect tenses due to the closeness of their meaning. See below for more. When people are talking about distant past, the assumption is that the activity discussed is recurrent (unless the context dictates otherwise), especially if we attach habitual time signal. For example: The notion of recent past is roughly limited to about last night. So, "tadi malam" means last night. Baru saja is roughly equivalent to "just now". However, it can also means within an hour or so. To express definite past, we can use the phrase "lalu", which roughly means ago. For example: Note that sometimes the phrase "yang lalu" is used instead of "lalu". They are equivalent. Literally, the word lalu means pass and yang means that or which. So, "dua jam yang lalu" literally means "[at] two hours that pass". Other words that may be useful to express ideas in the past: Future Tense. The same goes with future tense: We need to attach time signals. For example: To express undefinite future, Indonesian has these phrases: Note that the word akan must be placed right before the verb, just like sedang. See the example above. The difference between kelak and nanti is on the distance to the future they are. Their usage can be combined with the word akan. The first two sentences are equivalent, as are the last two. The difference is that the first two implies more distant future than the last two. How distant? It depends on the context. The good rule of thumb is kelak usually refers to a time frame of months or years in the future, whereas nanti refers to a much shorter time in the future than that (i.e. days). To specify definite future, the word akan can also be combined with other future time signals, such as besok ( = tomorrow). The word kelak or nanti can be combined with future time signals too in order to specify a definite future: Note that when specifying definite future, the word kelak and nanti is equivalent. So, "dua jam nanti" is equivalent as "dua jam kelak". Usually people still follow the "rule of thumb" above. So, "dua jam nanti" is used more often than "dua jam kelak"; and "dua tahun kelak" is used more often than "dua tahun nanti". However, people usually use "dua bulan nanti" and "dua bulan kelak" interchangeably. Example: Note that the word akan can also be combined with kelak or nanti. But kelak and nanti cannot be used together. For example: The two examples are equivalent. Another good time signal to use for future tense is "depan", which means next (in next week, etc.): Tense Combinations. As stated above, Indonesian has no complex tenses such as future perfect (i.e. "will have + infinitive"). The way Indonesian gets around with it is to throw in the appropriate time signals. This practice is influenced by Romance languages and sounds inherently unnatural to Indonesian people. However, you may do that and people may still be able to understand it, may be just a bit strange. Present Perfect. Due to English language influence, people began using the words sudah or telah to express present perfect. This is because Indonesian actually doesn't have present perfect. So: It's a bit weird, but it works. So you can use it if you want to. Future Perfect. We use the phrase "akan telah" or "akan sudah" to indicate future perfect. Conditional Perfect. Conditional perfect is expressed using "would have". Unfortunately, Indonesian has no way to express this. The usual translation is using "mungkin telah", but that would actually mean "may have". This is very awkward situation. It sounds very unnatural. The best way translating it is to translate the sentence in its entirety in a different way. 

Chapter 5. Plant Reproduction Laboratory ~ Flowers. &lt;br&gt; An orchid flower. This first laboratory exercise for Chapter 4 deals with the flowers of a ground orchid from Southeast Asia. The photograph on the right demonstrates the descriptive terminology that can be applied to this species. You may wish to read about to place this plant taxonomically and better understand unusual aspects of the structure of this flower. In reading the description below, be sure you understand how or why each bolded word applies to this specimen. Also, observe that the flowering-through-fruiting sequence is well demonstrated in the photograph because each flower is in a slightly different phase of its life cycle from bud to fruit. 4-1. Review the photograph of the inflorescence of the orchid. "Which one of these statements is true": &lt;br&gt; Following are a series of photographs of flowers from various plants. "Note that by clicking on the word" "Examine" "in each title, you can enlarge the particular photograph for closer examination". Read each question and the offered answers carefully. All parts of answer choice must be correct for that choice to be correct. 4-2. "The structure at B is": &lt;br&gt; 4-3. "Although the flowers in these three photographs appear very different, the following parts or floral structures are essentially the same": 4-4. "Which statement of the following applies to the structure indicated at E ": &lt;br&gt; « Return to Chapter 5 Answers to Chapter 4 Laboratory Questions: 

= Grammatical Introduction to Verbs = This introductory section may be a bit overwhelming, but is an overall look at verbs. The majority of this section will be covered in later chapters. Nevertheless, looking over this chapter may help you to familiarize yourself with verbs. Verbs are parts of speech which denote action. There are two main forms of verbs in Latin: • Principal Verbs (the main verb which is found in every sentence. e.g.,: vir ambulat = the man is walking) • Adjectival Verbs (also known as participles, gerunds and gerundives which describe the state of the described noun. e.g.,: vir ambulans = the walking man. The verb behaves as an adjective) Every sentence must have a verb. In a sense, the principal verb is the sentence and all the nouns, adverbs and participles are only describing the scenario of the verb. Thus in Latin this constitutes a sentence:  est. If you want to explain 'who' is or exists, you add a nominative substantive:  Cornēlia est. We now know Cornelia 'is'. But what is she? So we add an adjective.  Cornēlia est bona. Now we can see that Cornelia is good, but to elaborate further we can add an adverb:  Cornēlia vix est bona. Now we know that Cornelia is 'hardly' ("vix": hardly, scarcely, barely) good. Thus, in English, the shortest Latin sentence is: You are. in Latin:  es Examples. These two examples will demonstrate the difference between an adjectival verb and a principal verb. Personal Endings. Verbs in Latin are inflected to reflect the person who performs the action. English does the same to some extent in the verb to be: Latin, however, inflects all verbs, and is much more extensive than English, allowing writers and speakers of Latin to often drop the personal pronoun, as the performer of the action is understood by the formation of the verb. The Personal pronoun is only usually added for emphasis. In a way, the ending on Latin verbs are a type of pronoun. Moods. There are several moods. Each has its own uses to convey certain ideas. The most commons moods are: • Indicative • Subjunctive or Conjunctive • Imperative The two moods we will first learn are the imperative (commands and orders) and the indicative (declarative statements and factual questions). Voice. There are two constructions verbs can have regarding voice. Verbs can have either an active or passive voice. E.g. 'I smash the car.' 'smash' is an active verb construct. The passive is used when the nominative is affected by the verb. E.g. 'The car is smashed by me.' 'is smashed' is a passive construct. Tense. Tense in Latin comprises two parts: TIME and ASPECT. Time reflects when the action is occurring or did occur: past, present, or future. Aspect refers to the nature of the action: simple, completed, or repeated. The "completed" aspect is generally termed "perfective" and repeated aspect "imperfective." Theoretically, a verb could have nine tenses (combinations of time and aspect). However, Latin only has six, since some possible combinations are expressed by the same verb forms. Latin tenses do not correspond exactly to English ones. Below is a rough guide to tense in Latin. As is evident, some Latin tenses do "double duty." The Latin Present and Future Tenses can either express simple or progressive aspect. Particularly difficult to grasp is the Latin Perfect tense, which can either express an action completed from the point of view of the present ("I have just now finished walking"), or a simple action in past time (its "aorist" sense, from the old Indo European aorist tense, which Latin lost but is still present in Greek). Infinitive. The infinitive (impersonal) is the form of the verb which simply means 'to (verb)' e.g. 'to do', or 'to be', or 'to love', or 'to hate' etc. All forms which are not in the infinitive are in the finite (personalised) form. The infinitive has a -re at the end of the stem of the verb. The infinitive of 'to be' is an exception and is 'esse'. dēbeō currere nunc = I ought to run now. esse, aut nōn esse = To be, or not to be? Exercises. Answer these two questions about the infinitive and finite. Irregularities. Verbs which use the passive formation in an active sense are known as deponent. Verbs which don't have a form for every tense and mood are known as defective. You will meet a few words like this soon. Personal Pronouns. In case you do ever use a personal pronoun to emphasise the SUBJECT of the verb, you must remember that the personal pronoun must be in the nominative case and the number and person of the verb must match that of the subject. (Review Lesson 7 if unfamiliar with the terms person and subject). Principal Parts. When one looks up a verb in the dictionary, the principal parts are given. From these principal parts you can find the correct form of the verb for every circumstance. Exercises. Answer this question about principal parts. Using the Dictionary. All nouns are given in the nominative, as well as the declension and gender of the noun. Verbs are alphabetized using the 1st person singular (the first principal part) and the infinitive is given. Supplementary principal parts are given if the various other principal parts do not follow the standard pattern of formation from the infinitive and 1st person singular. =Verbs: Conjugation in the Present Imperfect = The present imperfect is the simplest tense. To form the present imperfect all that is required is to place the personal endings at the end of the verb stem. Thus, if you have the stem 'ama' (love), to make it 'I love' you place an ō at the end.  I love = amō (amaō*)  we love = amāmus Latin could add personal pronouns, however only for added emphasis and in conjunction with the corresponding person ending on the verb. Otherwise the sentence will not make sense. For example: ego amō = I (not you) love nōs amāmus = We (not you) love but that would be for special emphasis: It's I, not you, who loves. Here are the forms of the verb 'porta', carry, in the present imperfect tense:  portō I carry first person singular  portās thou carriest, you carry second person singular  portat he, she, it carries third person singular  portāmus we carry first person plural  portātis you (all) carry second person plural  portant they carry third person plural 'porto' can also be translated 'I am carrying' (present imperfect), 'I do carry' (present emphatic). 'I carry' is known as the 'present simple' tense in English. Again the 'a' gets dropped when the 'ō' is placed on porta. Porta, and ama are known as 1st conjugation verbs; in other words, verbs which have a stem ending in 'a'. There are three other conjugations, and below are some examples of verbs from each of the four conjugations (present imperfect tense): Each verb uses the same final letter or letters to indicate the 'subject' - I, thou, he/she/it, we, you, they. Before these final letters, the first conjugation has an 'a' (although when an 'o' is placed, the 'a' is often dropped), the second an 'e', and the third and fourth usually an 'i'. The third person plural forms in the third and fourth conjugations have a 'u'. These verb forms really should be learned by heart. The most common verb of all is irregular (see next lesson). Here is a table of the verb 'to be' in Latin, English, and four Romantic languages (French, Spanish, Italian and Portuguese) The personal endings are the same as in the four regular conjugations. Exercises. &lt;br&gt; Imperative Mood. The imperative mood conveys an order (e.g. Go!, Run!, Away Now!). The imperative mood is formed by simply using the stem of the verb. If the order is to a large group of people, or you are trying to show respect, you must use the -te suffix. amō eum = I love him. amā eum = Love him! amāte eum = Love (respectful, or plural) him! currō casam = I run home. curre casam = Run home! currite casam = Run (respectful, or plural) home! regō prudente = I rule wisely. rege prudente = Rule wisely! regite prudente = Rule (respectful order) wisely! 

External Links. This textbook can be heard read out aloud here: 

Eigenvalues and eigenvectors are related to fundamental properties of matrices. Motivations. Large matrices can be costly, in terms of computational time, to use, and may have to be iterated hundreds or thousands of times for a calculation. Additionally, the behavior of matrices would be hard to explore without important mathematical tools. One mathematical tool, which has applications not only for Linear Algebra but for differential equations, calculus, and many other areas, is the concept of "eigenvalues" and "eigenvectors". Eigenvalues and eigenvectors are based upon a common behavior in linear systems. Let's look at an example. Let and What happens with x and y if they are transformed by "A"? Well, But what is remarkable is that So when we operate on the vector x with the matrix "A", instead of getting a different vector (as we would normally do), we get the "same" vector x multiplied by some constant. And the same goes for vector y. We call the values 1 and -2 the "eigenvalues" of the matrix "A", and the vectors x and y are called "eigenvectors" for the matrix "A". Definitions. We now generalize this concept of when a matrix/vector product is the same as a product by a scalar as above: essentially if we have a "n"×"n" matrix A, we seek solutions in v to find the eigenvectors, and solutions in λ to find the eigenvalues for the equation How are we to do this? Let us rearrange the equation But (A-λI) is a matrix, so we are trying to solve Bv=0 where B=(A-λI), and this solution is merely the kernel of B, ker B. So the eigenvectors are in ker (A-λI), where λ is an eigenvalue. But how do we find the eigenvalues? Bv=0 has nonzero solution if |B| = det(B) is zero. So to find the eigenvalues, we let |A-λI|=0 and then solve for λ. We will thus obtain a polynomial equation over the complex numbers (eigenvalues can be complex), known as the "characteristic equation". The roots of the characteristic equation are the eigenvalues. Note that we exclude 0 as an eigenvector, because it is trivially a solution to Av=λv and is not really interesting to consider. Additionally, if the zero vector were to be included, it would allow for an infinite number of eigenvalues, since any value of λ satisfies A0=λ0. If we have an eigenvalue λ of a matrix A, together with a corresponding eigenvector, x, then any multiple of x is also an eigenvector for the same eigenvalue. To see that kx is also an eigenvector, follow this argument: If Ax=λx, then A(kx)=kAx=kλx=λ(kx). (Here k may be any scalar.) Thus, every multiple of an eigenvector is also an eigenvector. Note the asymmetry here: eigenvalues are unique, while an eigenvalue has many eigenvectors. &lt;/gallery&gt; &lt;/gallery&gt; &lt;/gallery&gt; Bold textÆə=== Finding eigenvalues and eigenvectors === Here are some examples of finding eigenvalues and eigenvectors using our definitions. Let Firstly, we expand |A-λI|=0 to find the eigenvalues: Now, elementary algebra tells us the roots of this equation are 3 and 2, and thus these are our eigenvalues. Now we can find our eigenvectors. Consider the first eigenvalue λ=3. To find our first eigenvector At this point we can row-reduce and back-substitute, but usually it suffices to guess the kernel since our matrix is small and we have linearly dependent columns. Now, observe: So, for any scalar a, the vector As noted above the eigenvalues of a matrix are uniquely determined, but for each eigenvalue there are many eigenvectors. We usually choose an eigenvector for some convenience such as "most whole number entries", "first entry is 1", or "length of the eigenvector is 1". Most Computer Algebra Systems choose unit vectors for eigenvectors. So here we may take formula_16 to be the eigenvector, for example. Similarly for our second eigenvalue λ=2, to find our second eigenvector: And so, our second eigenvector is chosen as Our eigenvalues then are λ=2,3, with eigenvectors formula_19, as may be checked by multiplying each by the given matrix. Problem set. Given the above, find the eigenvalues and eigenvectors of the following matrices (Answers follow to even-numbered questions): Applications. Eigenvalues and eigenvectors are not mere pretty facts about these vectors; they have relevant and important applications. Matrix powers. Let us first examine a certain class of matrices known as "diagonal" matrices: these are matrices in the form Now, observe that This is a useful property! However, the number of matrices to which we can apply this fact is clearly limited, so we ask ourselves whether we can transform a given matrix into a diagonal matrix. The answer to this question is "sometimes", but for the moment, we will only look at matrices for which this answer is "yes". What we seek is a matrix P such that where D is diagonal. If such a matrix P exists, we say that A is "diagonalizable". (Note that "xyx"-1 is often called a "similarity transformation"). Then by multiplying throughout forward by P-1, then by multiplying backward by P. Now, we have The PP-1 terms cancel to give We can calculate Dk easily, so we need to find P. It turns out (the entire proof is quite difficult) that we simply create a matrix from concatenating the linearly independent eigenvectors to create P. D, then, is the diagonal matrix containing the eigenvalues on the main diagonal corresponding to the associated eigenvectors (the eigenvalue in the first place corresponds to the eigenvector it is created from, in the first column). Example. Let's work through an example to show these ideas. So what do we do if we want to find A14? Let's use the method we've just described. Find the eigenvalues: Find the eigenvectors: The eigenvectors are then so put the eigenvectors together to form the matrix P Now -1 generated the eigenvector in the "first" column, and 4 generated the eigenvector in the "second" column, so form D in this way: We can easily calculate (-1)14=1, so we get and we have the fast method for creating inverses of 2×2 matrices: So now we can now directly multiply out Simplifying we get Problem set. Given the above, find the following matrix powers (Answers follow to even-numbered questions): Coupled ordinary differential equations. We can use the method of diagonalisation to solve coupled ordinary differential equations. For example, let x(t) and y(t) be differentiable functions and x' and y' their derivatives. The differential equations are relatively difficult to solve: but it has solution remembering this fact, we translate the ODEs into matrix form Diagonalise the square matrix, we get: we put then it follows that thus as discussed above the solutions are easy. We have for some constants C and D. Now that we get and so This method generalises well into higher dimensions. Coupled differential equations. Matrices, strangely enough, have a great use in relation to "calculus" in the calculation of solutions to coupled differential equations, where one differential equation has some function that depends on another differential equation. For example: Without going any further, the solution to these differential equations looks very difficult! However if we formulate this in terms of matrices, it becomes a little bit easier to analyze. Example. Let's take the above example, so Now form a vector: Then Now the problem becomes This is reminiscent of the differential equation we have already encountered in calculus, that of in which the solution is y = "c"ekt. We can make a wild guess then the solution to the above matrix equation will have a solution in a similar form. So let's try a solution v = weλt. Then D v = λweλt. Let us then try and substitute this guess solution into our equation: If we let we see that the equation above becomes, on dividing through by formula_61 (since it is never zero) But wait - this is the equation before to find the eigenvalues - and we have that the solution v = weλt is a solution if and only if λ is an eigenvalue of A and w is its corresponding eigenvector. The eigenvalues are 4, 2, with eigenvectors respectively. So we have two solutions and Note that if we have two solutions to the differential equation D v = Av, linear combinations of the two solutions will give the same solution. So then we have then the general solution: Separating into the first and second components we get our two solutions Problem set. Given the above solve the following problems (answers to even-numbered questions follow) Answers. Form the matrix The eigenvalues of this matrix are and the eigenvectors are So now and y("t") and x("t") can be read off by inspection. 

A vector space is a way of generalizing the concept of a set of vectors. For example, the complex number 2+3i can be considered a vector, since in some way it is the vector formula_1. The vector space is a "space" of such abstract objects, which we term "vectors". Some familiar friends. Currently in our study of vectors we have looked at vectors with real entries: formula_2, and so on. These are all vector spaces. The advantage we gain in abstracting to vector spaces is a way of talking about a space without any particular choice of objects (which define our vectors), operations (which act on our vectors), or coordinates (which identify our vectors in the space). Further results may be applied to more general spaces which may have infinite dimension, such as in Functional Analysis. Notations and concepts. We write a vector, like we have before, bold, but you should write these on paper underlined or with an arrow on top. So we write formula_3 for that vector. When we multiply a vector by a scalar number, we usually ascribe it a Greek letter, writing λv for the multiplication of v by a scalar λ. We write addition and subtraction of vectors as we have been doing before, x+y for the sum of vectors x and y. With scalar multiplication and adding vectors, we can move to our definition of a vector space. When we refer to an operation being 'closed' in a definition, we are saying that the result of the operation does not violate our definition. For example, if we are looking at the set of all integers, we can say that it is closed under addition, because adding any integers results in something inside the set of integers. However the set of integers is not closed under division, because dividing 3 by 2 (for example) doesn't result in a member of the set of integers. Definition. A "vector space" is a nonempty set V of objects, called "vectors", on which are defined two operations, called "vector addition" and "scalar multiplication", respectively, such that, for formula_4 and αformula_5, where F is a field x+y and αx are well-defined elements of V with the following properties: Alternative Definition. People who are familiar with group theory and field theory may find the following alternative definition more compact: Linear Spaces. The linear space is a very important vector space. Let n1, n2, n3, ..., nk be k elements of a field F. Then the ordered k-tuples (n1, n2, n3, ..., nk) form a vector space with addition being the sum of the corresponding numbers, and scalar multiplication by an element of F being the result of multiplying each one in the k-tuple. This would then be the k dimensional linear space. Subspaces. A "subspace" is a vector space inside a vector space. When we look at various vector spaces, it is often useful to examine their subspaces. The subspace S of a vector space V means that S is a sub"set" of V and that it has the following key characteristics Any subset with these characteristics is a vector space. The trivial subspace. The singleton set with the zero vector ({0}) is a subspace of every vector space. Scalar multiplication closure: "a" 0=0 for all "a" in R Addition closure: 0+0=0. Since 0 is the only member of the set we need to check only 0 Zero vector: 0 is the only member of the set and it is the zero vector. Examples. Let us examine some subspaces of some familiar vector spaces, and see how we can prove that a certain subset of a vector space is in fact a subspace. A slightly less trivial subspace. In R2, the set V of all vectors from R2 of the form (0,α) where α is in R is a subspace Scalar multiplication closure: "a" (0,α) = (0,a α) and a α is in R Addition closure: (0,α) +(0,β) =(0, α + β) and α + β is in R Zero vector: taking α to be zero in our definition of (0, α) in V we get the zero vector (0,0) A whole family of subspaces. Pick any number from R, say ρ. Then the set V of all vectors of the form (α, ρα) is a subspace of R2 Scalar multiplication closure: "a" (α, ρα) = (aα, ρaα) which is in V. Addition closure: (α, ρα) +(β, ρβ) =(α + β, ρα + ρβ) = (α+β, ρ(α+β)) which is in V Zero vector: taking α to be zero in our definition we get (0, ρ0) = (0,0) in V. That means V2 = the set of all vectors of the form (α,2α) is a subspace of R2 and V3 = the set of all vectors of the form (α,3α) is a subspace of R2 and V4 = the set of all vectors of the form (α,4α) is a subspace of R2 and V5 = the set of all vectors of the form (α,5α) is a subspace of R2 and Vπ = the set of all vectors of the form (α,πα) is a subspace of R2 and V√2 = the set of all vectors of the form formula_11 is a subspace of R2 As you can see, even a simple vector space like R2 can have many different subspaces. Linear Combinations, Spans and Spanning Sets, Linear Dependence, and Linear Independence. Linear Combinations. Definition: Assume formula_12 is a "vector space" over a field formula_13 and formula_14 is a nonempty subset of formula_12. Then a vector formula_16 is said to be a linear combination of elements of formula_14 if there exists a finite number of elements formula_18 and formula_19 such that formula_20. Spans. Definition: Assume formula_12 is a "vector space" over a field formula_13. The set of all linear combinations of formula_23 is called the span of formula_24. This is sometimes denoted by formula_25. Note that formula_26 is a subspace of formula_12. Proof: Consider closure under addition and scalar multiplication for two vectors, x and y, in the span of the vectors formula_28 formula_29 formula_30 formula_31, which is also contained in the set. formula_32, which is also contained in the set. Spanning Sets. Definition: Assume formula_12 is a "vector space" over a field formula_13 and formula_24 are vectors in such a vector space. The set formula_36 is a spanning set for the vector space formula_12 if and only if every vector in formula_12 is a linear combination of formula_24. Alternately, formula_40 Linear Independence. Definition: Assume formula_12 is a "vector space" over a field formula_13 and formula_43 is a finite subset of formula_12. Then we say formula_14 is linearly independent if formula_46 implies formula_47. Linear independence is a very important topic in Linear Algebra. The definition implies that linearly dependent vectors may form the nulvector as a non-trivial combination, from which we may conclude that one of the vectors can be expressed as a linear combination of the others. If we have a vector space V spanned by 3 vectors formula_48 we say that v1, v2, and v3 are linearly dependent if there is a combination of one or two of them that can produce a third. For instance, if one of the following equations: can be satisfied, then the vectors in V are said to be linearly dependant. How can we test for linear independence? The definition sets it out to us: If V is a vector space spanned by 3 vectors of length N: and we try to test whether these 3 vectors are linearly independent, we form the equations: and solve them. If the only solution is then the 3 vectors are linearly independent. If there is another solution they are linearly dependent. ?????? We can say that for V to be linearly independent it must satisfy this condition: Where we are using 0 to denote the null vector in V. If formula_56 is square and invertable, we can solve this equation directly: And if we know that formula_59 is zero, then we know that the system is linearly independent. If, however, formula_56 is not square, or if it is not invertable, we can try the following technique: Multiply through by the transpose matrix: Find the inverse of formula_62, and multiply through by the inverse: Cancel the terms: And our conclusion: This again means that V is linearly independent. Span. A span is the set of all possible vectors that are in a given vector space. Basis. A basis for a vector space is the least amount of linearly independent vectors that can be used to describe the vector space completely. The most common basis vectors are the kronecker vectors, also called canonical basis: In the cartesian graphing space, we say an ordered triple of coordinates is defined as: And we can make any point (x, y, z) by combining the kronecker basis vectors: Some theorems: Bases and Dimension. If a vector space V is such that:&lt;br&gt; it contains a linearly independent set B of N vectors, and&lt;/br&gt;&lt;br&gt; any set of N + 1 or more vectors in V is linearly dependent,&lt;/br&gt; then V is said to have dimension N, and B is said to be a basis of V. TODO. Tell about what is a "basis" in a vector space and about coordinate transformations. (this article contains an abstract definition of a "basis" which is a generalization of a basis in vector space and can be used as the foundation to explain about bases and coordinate transformations.) Discuss the geometry of subspaces (points, lines, planes, hypersurfaces) and connect them to the geometry of solutions of linear systems. Connect the algebra of subspaces and linear combinations of vectors to the algebra of linear systems. 

 Setup and Installation. Note: Before contributing, check out the discussion page. How to write your examples. Learning the Language. Advanced PHP. Security. "See also:" External links.  __NOEDITSECTION__ 

Appendices.  __NOEDITSECTION__ 



Kinematics is the description of motion. The motion of a point particle is fully described using three terms - position, velocity, and acceleration. For real objects (which are not mathematical points), "translational kinematics" describes the motion of an object's center of mass through space, while "angular kinematics" describes how an object rotates about its centre of mass. In this section, we focus only on translational kinematics. Position, displacement, velocity, and acceleration are defined as follows. Position. "Position" is a relative term that describes the location of an object RELATIVE to some chosen stationary point that is usually described as the "origin". A vector is a quantity that has both magnitude and direction, typically written as a column of scalars. That is, a number that has a direction assigned to it. In physics, a vector often describes the motion of an object. For example, Warty the Woodchuck goes 10 meters towards a hole in the ground. We can divide vectors into parts called "components", of which the vector is a sum. For example, a two-dimensional vector is divided into x and y components. Displacement. Displacement answers the question, "Has the object moved?" Note the formula_1 symbol. This symbol is a sort of "super equals" symbol, indicating that not only does formula_2 EQUAL the displacement formula_3, but more importantly displacement is operationally defined by formula_2. We say that formula_2 operationally defines displacement, because formula_2 gives a step by step procedure for determining displacement. Namely: Be sure to note that displacement is not the same as distance travelled. For example, imagine travelling one time along the circumference of a circle. If you end where you started, your displacement is zero, even though you have clearly travelled some distance. In fact, displacement is an average distance travelled. On your trip along the circle, your north and south motion averaged out, as did your east and west motion. Clearly we are losing some important information. The key to regaining this information is to use smaller displacement intervals. For example, instead of calculating your displacement for your trip along the circle in one large step, consider dividing the circle into 16 equal segments. Calculate the distance you travelled along each of these segments, and then add all your results together. Now your total travelled distance is not zero, but something approximating the circumference of the circle. Is your approximation good enough? Ultimately, that depends on the level of accuracy you need in a particular application, but luckily you can always use finer resolution. For example, we could break your trip into 32 equal segments for a better approximation. Returning to your trip around the circle, you know the true distance is simply the circumference of the circle. The problem is that we often face a practical limitation for determining the true distance travelled. (The travelled path may have too many twists and turns, for example.) Luckily, we can always determine displacement, and by carefully choosing small enough displacement steps, we can use displacement to obtain a pretty good approximation for the true distance travelled. (The mathematics of calculus provides a formal methodology for estimating a "true value" through the use of successively better approximations.) In the rest of this discussion, I will replace formula_7 with formula_8 to indicate that small enough displacement steps have been used to provide a good enough approximation for the true distance travelled. Velocity. [Δ, "delta", upper-case Greek D, is a "prefix" conventionally used to denote a "difference".] Velocity answers the question "Is the object moving now, and if so - how quickly?" Once again we have an "operational" definition: we are told what steps to follow to calculate velocity. Note that this is a definition for average velocity. The displacement Δ"x" is the vector sum of the smaller displacements which it contains, and some of these may subtract out. By contrast, the distance travelled is the scalar sum of the smaller distances, all of which are non-negative (they are the "magnitudes" of the displacements). Thus the distance travelled can be larger than the magnitude of the displacement, as in the example of travel on a circle, above. Consequently, the average velocity may be small (or zero, or negative) while the speed is positive. If we are careful to use very small displacement steps, so that they come pretty close to approximating the true distance travelled, then we can write the definition for instantaneous velocity as [δ is the lower-case "delta".] Or with the idea of limits from calculus, we have: [d, like Δ and δ, is merely a "prefix"; however, its use definitely specifies that this is a sufficiently small difference so that the error--due to stepping (instead of smoothly changing) the quantity--becomes negligible.] Acceleration. Acceleration answers the question "Is the object's velocity changing, and if so - how quickly?" Once again we have an operational definition. We are told what steps to follow to calculate acceleration. Again, also note that technically we have a definition for average acceleration. As for displacement, if we are careful to use a series of small velocity changes, then we can write the definition for instantaneous acceleration as: Or with the help of calculus, we have: Vectors. Notice that the definitions given above for displacement, velocity and acceleration included little arrows over many of the terms. The little arrow reminds us that direction is an important part of displacement, velocity, and acceleration. These quantities are vectors. By convention, the little arrow always points right when placed over a letter. So for example, formula_9 just reminds us that velocity is a vector, and does not imply that this particular velocity is rightward. Why do we need vectors? As a simple example, consider velocity. It is not enough to know how fast one is moving. We also need to know which direction we are moving. Less trivially, consider how many different ways an object could be experiencing an acceleration (a change in its velocity). Ultimately, there are three distinct ways an object could accelerate: More general accelerations are simply combinations of 1 and 3 or 2 and 3. Importantly, a change in the direction of motion is just as much an acceleration as is speeding up or slowing down. In classical mechanics, no direction is associated with time (you cannot point to next Tuesday). So the definition of formula_10 tells us that acceleration will point wherever the change in velocity formula_11 points. Understanding that the direction of formula_12 determines the direction of formula_13 leads to three non-mathematical but very powerful rules of thumb: Again, more general motion is simply a combination of 1 and 3 or 2 and 3. Using these three simple rules will dramatically help your intuition of what is happening in a particular problem. In fact, much of the first semester of college physics is simply the application of these three rules in different formats. = Equations of motion (constant acceleration) = A particle is said to move with constant acceleration if its velocity changes by equal amounts in equal intervals of time, no matter how small the intervals may be Since acceleration is a vector, constant acceleration means that both direction and magnitude of this vector don't change during the motion. This means that average and instantaneous acceleration are equal. We can use that to derive an equation for velocity as a function of time by integrating the constant acceleration. Giving the following equation for velocity as a function of time. To derive the equation for position we simply integrate the equation for velocity. Integrating again gives the equation for position. The following are the equations of motion: The following equations can be derived from the two equations above by combining them and eliminating variables. What does force in motion mean? Force means strength and power. Motion means movement. That’s why we need forces and motions in our life. We need calculation when we want to know how fast things go, travel and other things which have force and motion. How do we calculate the speed? If you want to calculate the average speed, distance travelled or time taken you need to use this formula and remember it:  formula_14 This is an easy formula to use, you can find the distance travelled, time taken or average speed, you need at least 2 values to find the whole answer. Is velocity the same thing as speed? Velocity is a vector quantity that refers to "the rate at which an object changes its position", whereas speed is a scalar quantity, which cannot be negative. Imagine a kid moving rapidly, one step forward and one step back, always returning to the original starting position. While this might result in a frenzy activity, it would result in a zero velocity, because the kid always returns to the original position, the motion would never result in a change in position, in other words formula_15 would be zero. Speed is measured in the same physical units of measurement as velocity, but does not contain an element of direction. Speed is thus the magnitude component of velocity. Velocity contains both the magnitude and direction components. You can think of velocity as the displacement/duration, whereas speed can be though as distance/duration. Acceleration. When a car is speeding up we say that it is accelerating, when it slows down we say it is decelerating. How do we calculate it? When we want to calculate it, the method goes like that: A lorry driver brakes hard, and slows from 25 m/s to 5 m/s in 5 seconds. What was the vehicle's acceleration?  formula_16 What is initial velocity and final velocity? Initial velocity is the beginning before motion starts or in the middle of the motion, final velocity is when the motion stops. There is another way to calculate it and it is like that This equations which are written is the primary ones, which means that when you don’t have lets say final velocity, how will you calculate the equation? This is the way you are going to calculate. Observing motion. When you want to know how fast an athletic person is running, what you need is a stopwatch in your hand, then when the person starts to run, you start the stopwatch and when the person who is sprinting stops at the end point, you stop the watch and see how fast he ran, and if you want to see if the athlete is wasting his energy, while he is running look at his movement, and you will know by that if he is wasting his energy or not. This athletic person is running, and while he is running the scientist could know if he was wasting his energy if they want by the stop watch and looking at his momentum. Measuring acceleration. Take a slope, a trolley, some tapes and a stop watch, then put the tapes on the slope and take the trolley on the slope, and the stopwatch in your hand, as soon as you release the trolley, start timing the trolley at how fast it will move, when the trolley stops at the end then stop the timing. After wards, after seeing the timing , record it, then you let the slope a little bit high, and you will see, how little by little it will decelerate. = Newton = Isaac Newton was an English physicist, mathematician (described in his own day as a "natural philosopher") , astronomer and alchemist. Newton is one of the most influential scientists of all time, and he is known, among other things, for contributing to development of classical mechanics and for inventing, independently from "Gottfried Leibniz", calculus. Newton's laws of motion. Newton is also known by his three laws of motion, which describe the relationship between a "body" and the "forces acting upon it", and "its motion in response to said forces". =Symbols= Some useful symbols seen and that we will see: Force. A force is any interaction that tends to change the motion of an object. In other words, a force can cause an object with mass to change its velocity. Force can also be described by intuitive concepts such as a push or a pull. A force has both magnitude and direction, making it a vector quantity. It is measured in the SI unit of "newtons" and represented by the symbol formula_17. How to calculate the force? When we want to calculate the force, and we have the "mass" and "acceleration", we can simply use the simple formula stated in the Newton's second law above, that is formula_20, where formula_18 is the mass (or the amount of matter in a body), and formula_19 is the acceleration. Note that the Newton’s second law is defined as a numerical measure of inertia. What is inertia? Inertia is the tendency of a body to maintain its state of rest or uniform motion, unless acted upon by an external force. Robert Hooke. Robert Hooke was an English polymath who played an important role in the scientific revolution, through both experimental and theoretical work. Hooke's law. Hooke's law is a principle of physics that states that the force formula_17 needed to extend or compress a spring by some distance formula_26 is proportional to that distance, or algebraically formula_27, where formula_28 is a constant factor characteristic of the spring, its stiffness. 

 Lektion 4 Lesestück 4-1 ~ Eine Geschichte über Zürich. Although this short story contains quite a number of impressive German nouns and adjectives, with the aid of "Vokabeln 4-1" following you should have no trouble reading and understanding it. The passage makes considerable use of the German genitive case (English possessive case), which you have not yet learned. However, a clue applicable here: translate "des" as "of the" or "of" and note there are other "der"-words that also mean "of the". Vokabeln 4-1.  die Alpen Alps  der Ausfluss outlet, effluence (of a lake)  die Bankinstitute banking institutes  die Bankenwirtschaft banking business  das Ende end  die Großbanken major banks  die Hauptstadt capital city  das Haus house  der Kanton canton (Swiss state)  das Lesestück reading passage  die Schweiz Switzerland  die Sicht view  der Sitz office  das Wetter weather  das Zentrum center (centre)  das Zürich Zurich (city and canton in Switzerland)  der Zürichsee Lake Zurich  d.h. (das heißt) i.e. ("that is" in Latin)  Glarner Alpen Glarner Alps  man hat... one has...  nach Hause (toward) home (compare: "zu Hause" = "at home")  anrufen call, telephone  geben (gab, gegeben) give  kommen (kam, gekommen) come  liegen (lag, gelegen) lie (lay, lain)  am (an dem) at the  ausgesprochen markedly  bei in  beiden two  etliche a number of, quite a few, several  gleichnamig same named  größte largest  klar clear  klein small  neben besides  nördlich northern  schweizer of or pertaining to Swiss Grammatik 4-1 ~ Introduction to adjectives. An is a part of speech which can be thought of as a "describing word"—typically, an adjective modifies a noun. In both English and German, adjectives come before the noun they describe or modify. In many other languages (such as French) they usually come after the noun. Here are some examples of adjectives (underlined) you have already encountered: Because nouns are capitalized in German, it is fairly obvious in these sentences where the adjectives occur: just before the nouns they modify. Note how the endings on German adjectives can change, depending upon the noun ("keinen Käse"; "klarem Wetter"; "gute Sicht")—specifically, the gender and case of the noun they are modifying. Before explaining the basic rules governing adjective endings, you need to have a better understanding of person, gender, and case in German nouns—concepts that will be explored in the next few lessons. Finally, realize that the ordinal numbers you learned in Lektion 3 are, in fact, adjectives—subject to the same rules governing word endings for adjectives. Gespräch 4-1 ~ Das neue Mädchen. This short conversational passage contains more examples of adjectives. Vokabeln 4-2.  die Brünette brunette  die Haare hair(s)  das Mädchen girl  das Ferkel piglet  gefallen appeal to  glauben believe  heißen name, call  mag like, desire, wish  dort there  (dort) drüben over there  dunkel dark  ihr her  hübsch cute  klein short  lang long  neue new  wenn if  wer? who? Grammatik 4-2 ~ Nouns and pronouns in the accusative and dative. As was noted previously when the concept of case was introduced for pronouns (Grammatik 2-2), there are four cases used in German. Recall that the nominative case in German corresponds to the "subjective case" in English and applies to nouns and pronouns used in a sentence as the subject of a verb. Nouns (and pronouns) that are used as objects of transitive (action) verbs are in the English objective case. If these are direct objects (recipients of the action of a verb), then these nouns are in the accusative case in German. If indirect objects, then these nouns are in the dative case in German. Essentially, the English "objective case" is divided, in German, into an accusative case used for direct objects and a dative case used for indirect objects. Pronouns. For comparison with English, recall that the singular personal pronouns ("nominative case") are "I", "you", and "he/she/it" (1st, 2nd, and 3rd persons). The "objective case", personal pronouns in English are "me", "you", and "him/her/it"—and are used for both direct and indirect objects of verbs. For example: The German accusative case, personal pronouns (singular) are: "mich, dich, ihn/sie/es". The German dative case, personal pronouns (singular) are: "mir", "dir", "ihm/ihr/ihm". Thus, the above English example sentence becomes, in German: Because "mir" is a dative pronoun, there is no need in German to use a modifier as in English, where "to" is used as a signal of an indirect object. The following table summarizes the German pronouns in three cases for both singular and plural number: Recall from Gespräch 2-1 the "incomplete" sentence "Und Ihnen?" ('And you?'). Note that the pronoun agrees in case (here, dative) with the implied sentence — "Und wie geht es Ihnen?" The same rule is evident in Gespräch 1-1 ("Und dir?"). Such agreement is important to convey the correct meaning. Tables giving the German personal pronouns in all cases can be found in an appendix: Pronoun Tables. Nouns. Nouns do not change their form (spelling) relative to case in German; instead, a preceding article indicates case. You have learned the nominative case definite and indefinite articles (Grammatik 3-3: "der", "die", "das" and "ein", "eine". "ein") for each of the three noun genders. Now we will learn the "accusative" (used to signal a direct object) and "dative" (used to signal an indirect object) articles. First, the definite articles: This table might seem a bit overwhelming (and there is yet one more case in German: the genitive!), but some points to note can make memorizing much easier. First, as you can see from the table, "gender" does not really exist for plural nouns. No matter what the noun gender in its singular number, its plural always has the same set of definite articles: "die", "die", "den" for nominative, accusative, and dative cases. The plural "der"-words are similar to the feminine singular "der"-words, differing only in the dative case. Another point: the dative for both masculine and neuter nouns is the same: "dem". Finally, for feminine, neuter, and plural nouns, there is no change between nominative and accusative cases. Thus, only for masculine nouns is there a definite article change in the accusative compared with the nominative. The following examples demonstrate the use of the definite article in various parts of speech: In the last example, you need to know that in both English and German, the noun (or pronoun) that follows the verb 'to be' is a predicate noun, for which the correct case is the nominative. That is why, in English, 'It is I' is grammatically correct and 'It is me' is simply incorrect. The indefinite articles are as follows: Of course, there are no plural indefinite articles in German or English ("ein" means "a". "an", or "one"). It is important to see that there is a pattern in the case endings added to "ein" related to the "der"-words in the definite articles table above. For example, the dative definite article for masculine nouns is "dem"—the indefinite article is formed by adding "-em" onto "ein" to get "einem". The dative definite article for feminine nouns is "der"—the indefinite is "ein" plus "-er" or "einer". These ending changes will be covered in greater detail in a future lesson. You will see that there are a number of words (adjectives, for example) whose form relative changes by addition of these endings to signal the case of the noun they modify. Finally, we can see a pattern relationship between these "endings" and the 3rd person pronouns as well: We could construct a similar table to compare the definite articles to the 3rd person pronouns. And in that case, we would also see how the plural definite articles ("die", "die", "den") compare with the third person plural pronouns ("sie", "sie", "ihnen"). Grammatik 4-3 ~ Interrogatives. You have encountered nearly all of the interrogatives commonly used in German (review Grammatik 1-2):  "wann" when  "warum" why "Warum sind Sie müde?"  "was" what "Was ist das?"  "wer" who "Wer ist das Mädchen?"  "wie" how "Wie geht es dir?"  "wieviel" how much "Wieviel Uhr ist es?"  "wo" where "Wo ist das Buch?"  "wohin" where (to) "Wohin gehst du?" In a question, interrogatives replace the unknown object and establish the class of answer expected. Note that the English construction for some of the questions differs from the German in that the former uses the progressive form of "do". Übersetzung 4-1. Translate the following sentences into German: 

=About the "Common uses in Physics"= While these are indeed common usages, it should be pointed out that there are many other usages and that other letters are used for the same purpose. The reason is quite simple: there are only so many symbols in the Greek and Latin alphabets, and scientists and mathematicians generally do not use symbols from other languages. It is a common trap to associate a symbol exclusively with some particular meaning, rather than learning and understanding the physics and relations behind it. =See Also= Greek alphabet on Wikipedia 

This book is a guide to Cascading Style Sheets (CSS), a technique widely used in web pages including Wikipedia to describe their visual style and appearance. CSS can take HTML to new places creatively and functionally. Once you learn how to style mark-up, you can additionally learn JavaScript functions that make dynamic web pages. 

Lektion Fünf &lt;br&gt; Gespräch 5-2 ~ Der Engländer in Österreich. &lt;br&gt; Wenn er auf den Kontinent fährt, wandert Herr Standish gern. Heute früh fährt er in die Stadt St. Pölten in Niederösterreich. Er spricht mit einer fremden Frau: Vokabeln 5A.  das Abendessen supper (evening meal)  [das] Österreich Austria  die Ecke corner  das Frühstück breakfast  das Hotel hotel  der Kilometer kilometer  die Küche cooking, cuisine  der Kontinent continent (Europe)  [das] Niederösterreich (federal state of) Lower Austria  das Rathaus city hall  das Restaurant restaurant  die Stadt city  Bitte sehr You're welcome  Entschuldigen Sie Pardon me, excuse me  Es gibt dort... There is there...  Gibt es...? Is there..?  Guten Tag good day (parting)  immer geradeaus straight on ahead  können Sie could you (polite form)  Wie bitte? Pardon me? (polite "come again?")  empfehlen recommend  fahren travel  kommen come, go, get  wandern wander  sagen say, tell  sprechen speak  anderer, andere, anderes other  besonders especially  bitte please  das that  dann then  darin therein  ein a (indefinite article)  eins one (cardinal number)  fremd unknown  gern gladly  gleich just, right (correct), right here, same  heute früh this morning  hier here (in this place)  ich I (personal pronoun)  links left (direction)  neben next to  rechts right (direction)  ungefähr approximately  von of ("Rathaus von St. Pölten" = St. Polten City Hall)  wie how (interrogative)  wo where (interrogative)  zu to ("zum" = contraction of "zu dem") Andere Wörter 4A.  der Bahnhof train station  der Flughafen airport  die Polizeiwache police station  die Post post office  genau exact(ly)  heute today Lesestück 5-1 ~ Eine Geschichte über St. Pölten. Niederösterreich ist sowohl flächenmäßig als auch nach Einwohnern das größte der neun österreichischen Bundesländer. Sankt Pölten ist die Landeshauptstadt von Niederösterreich. Der Name St. Pölten geht auf den heiligen Hippolytos zurück, nach dem die Stadt benannt wurde. Die Altstadt befindet sich dort, wo vom 2. bis zum 4. Jahrhundert die Römerstadt "Aelium Cetium" stand. 799 wurde der Ort als "Treisma" erwähnt. Das Marktrecht erhielt St. Pölten um 1050, zur Stadt erhoben wurde es 1159. Bis 1494 stand St. Pölten im Besitz des Bistums Passau, dann wurde es landesfürstliches Eigentum. Bereits 771 findet sich ein Benediktinerkloster, ab 1081 gab es Augustiner-Chorherren, 1784 wurde deren Kollegiatsstift aufgehoben, das Gebäude dient seit 1785 als Bischofssitz. Zur Landeshauptstadt von Niederösterreich wurde St. Pölten mit Landtagsbeschluss vom 10. Juli 1986, seit 1997 ist es Sitz der Niederösterreichischen Landesregierung. &lt;br&gt; Vokabeln 5B.  Die Altstadt old town  Der Augustiner Augustinian  Der Besitz possession, holding  Das Bistum diocese  Der Bischofssitz bishop's see (a seat of a bishop's authority)  Die Bundesländer federal states  Die Chorherren men's choir  Das Eigentum proprietorship  Die Einwohner inhabitants  Das Gebäude premises  Die Geschichte history  Das Jahrhundert century  Das Kloster monastery, friary  Das Kollegiatsstift monastery college  Die Landeshauptstadt regional or state capital city  Die Landesregierung provincial (state) government  Der Landtagsbeschluss day of jurisdictional reorganization  Das Marktrecht right to hold markets  Der Name name  Der Ort place, spot, city  Die Römerstadt Roman town  Der Sitz official place  Bistum Passau a dioecian region in Bavaria  sowohl... als auch both... and  zurück auf goes back to  aufheben (hob auf, aufgehoben) merged in (or turned into?)  befinden sich situated, located  (befand sich, haben sich befunden)  finden sich* found (located)  benennen (benannte, benannt) call (as to label)  erhalten (erhielt, erhalten) receive  erheben (erhob, erhoben) arise, raise  erwähnen (erwähnte, erwähnt) mention  stehen (stand, gestanden) stand (stood, stood)  werden (wurde, [ist]geworden) become  ab from  auf up  bereits already  bis until, by, up to  flächenmäßig (no direct translation) ~ when measured in surface  heilig holy  landesfürstlich baronial or princely (holdings)  nach in terms of  um around 

Introduction. Mathematicians have been, for the past five hundred years or so, obsessed with proofs. They want to prove everything, and in the process proved that they can't prove everything (see this). This chapter will introduce the axiomatic approach to mathematics, and several types of proofs. Direct proof. The direct proof is relatively simple — by logically applying previous knowledge, we "directly prove" what is required. Example 1 Prove that the sum of any two even integers formula_1 and formula_2 is even. Solution 1 We know that since formula_1 and formula_2 are even, they must have 2 as a factor. Then, we can write the following: Then: by the distributive property of integers The number formula_9 clearly has 2 as a factor, which implies it is even. Therefore, formula_10 is even. Example 2 Prove the following statement for non-zero integers formula_11: If formula_12 divides formula_13 and formula_13 divides formula_15 , then formula_12 divides formula_15 . Solution 2 If an integer formula_1 divides an integer formula_2 , then we can write formula_20 , for some non-zero integer formula_21 . So let's say that formula_22 and formula_23 , for some non-zero integers formula_21 and formula_25 . Then: by the associative property of integer multiplication. But since formula_21 and formula_25 are integers, their product formula_29 must also be an integer. Therefore, formula_15 is the product of some integer multiplied by formula_12 , so we get that formula_12 divides formula_15 . Mathematical induction. Deductive reasoning is the process of reaching a conclusion that is guaranteed to follow. For example, if we know then we can conclude: Induction is the opposite of deduction. To induce, we observe how things behave in specific cases and from that we draw conclusions as to how things behave in the general case. Suppose we want to show that a statement (let us call it formula_34 for easier notation) is true for all natural numbers. This is how induction a proof by induction works: To understand how the last step works, notice the following Example 1 Show that the identity holds for all positive integers. Solution Firstly, we show that it holds for 1 Suppose the identity holds for some natural number "k": This supposition is known as the induction hypothesis. We assume it is true, and aim to show that, is also true. We proceed which is what we have set out to show. Since the identity holds for 3, it also holds for 4, and since it holds for 4 it also holds for 5, and 6, and 7, and so on. There are two types of mathematical induction: strong and weak. In weak induction, you assume the identity holds for certain value k, and prove it for k+1. In strong induction, the identity must be true for any value lesser or equal to k, and then prove it for k+1. Example 2 Show that n! &gt; 2n for n ≥ 4. Solution The claim is true for n = 4. As 4! &gt; 24, i.e. 24 &gt; 16. Now suppose it's true for n = k, k ≥ 4, i.e. it follows that We have shown that if for n = k then it's also true for n = k + 1. Since it's true for n = 4, it's true for n = 5, 6, 7, 8 and so on for all n. Example 3 Show that Solution Suppose it's true for n = k, i.e. it follows that We have shown that if it's true for n = k then it's also true for n = k + 1. Now it's true for n = 1 (clear). Therefore it's true for all integers. Exercises. 1. Prove that formula_53 2. Prove that for n ≥ 1, where xn and yn are integers. 3. Note that Prove that there exists an explicit formula for formula_57 4. The sum of all of the interior angles of a triangle is formula_58; the sum of all the angles of a rectangle is formula_59. Prove that the sum of all the angles of a polygon with "n" sides, is formula_60. Proof by contradiction. The idea of a proof by contradiction is to: √2 is irrational. As an example, we shall prove that formula_61 is not a rational number. Recall that a rational number is a number which can be expressed in the form of p/q, where p and q are integers and q does not equal 0 (see the 'categorizing numbers' section here). First, assume that formula_61 is "rational":&lt;br&gt; where "a" and "b" are coprime (i.e. integers with no common factors, with greatest common divisor 1). If "a" and "b" are not coprime, we remove all common factors. In other words, "a/b" is in simplest form. Now, continuing: We have now found that "a2" is some integer multiplied by 2. Therefore, "a2" must be divisible by two. If "a2" is even, then "a" must also be even, for an odd number squared yields an odd number. Therefore we can write "a = 2c", where "c" is another integer. We have discovered that "b2" is also an integer multiplied by two. It follows that "b" must be even. We have a contradiction! Both "a" and "b" are even integers. In other words, both have the common factor of 2. But we already said that "a/b" is in simplest form, with no common factors. Since such a contradiction has been established, we must conclude that our original assumption was false. Therefore, √2 is irrational. Contrapositive. Some propositions that take the form of "if xxx then yyy" can be hard to prove. It is sometimes useful to consider the "contrapositive" of the statement. Before I explain what contrapositive is let us see an example is harder to prove than although they mean the same thing. So instead of proving the first proposition directly, we prove the second proposition instead. If "A" and "B" are two propositions, and we aim to prove we may prove the equivalent statement instead. This technique is called proof by contrapositive. To see why those two statements are equivalent, we show the following boolean algebra expressions is true (see Logic) (to be done by the reader). Exercises. 1. Prove that there is no perfect square number for 11,111,1111,11111... 2. Prove that there are infinitely number of "k"'s such that, 4"k" + 3, is prime. (Hint: consider N = p1p2...pm + 3) Reading higher mathematics. This is some basic information to help with reading other higher mathematical literature. "... to be expanded" Quantifiers. Sometimes we need propositions that involve some description of rough quantity, e.g. "For "all" odd integers x, x2 is also odd". The word "all" is a description of quantity. The word "some" is also used to describe quantity. Two special symbols are used to describe the quanties "all" and "some" Example 1&lt;br&gt; The proposition: can be expressed symbolically as: Example 2&lt;br&gt; The proposition: can be expressed symbolically as: This proposition is false. Example 3&lt;br&gt; Consider the proposition concerning (z = x'y' + xy): can be expressed symbolically as: This proposition is true. Note that the order of the quantifiers is important. While the above statement is true, the statement is false. It asserts that there is one value of y which is the same for all x for which z=1. The first statement only asserts that there is a y for each x, but different values of x may have different values of y. Negation. Negation is just a fancy word for the opposite, e.g. The "negation" of "All named Britney can sing" is "Some named Britney can't sing". What this says is that to disprove that all people named Britney can sing, we only need to find one named Britney who can't sing. To express symbolically: Similarly, to disprove we only need to find one odd number that doesn't satisfy the condition. Three is odd, but 3×3 = 9 is also odd, therefore the proposition is FALSE and is TRUE In summary, to obtain the "negation" of a proposition involving a quantifier, you replace the quantifier by its opposite (e.g. formula_67 with formula_77) and the "quantified proposition" (e.g. "x is even") by its negation (e.g. "x is odd"). Example 1 is a true statement. Its negation is Axioms and Inference. If today's mathematicians were to describe the greatest achievement in mathematics in the 20th century in one word, that word will be abstraction. True to its name, abstraction is a very abstract concept (see Abstraction). In this chapter we shall discuss the "essence" of some of the number systems we are familiar with. For example, the real numbers and the rational numbers. We look at the most fundamental properties that, in some sense, "define" those number systems. We begin our discussion by looking at some of the more obscure results we were told to be true Most people simply accept that they are true (and they are), but the two results above are simple consequences of what we believe to be true in a number system like the real numbers! To understand this we introduce the idea of axiomatic mathematics (mathematics with simple assumptions). An axiom is a statement about a number system that we assume to be true. Each number system has a few axioms, from these axioms we can draw conclusions (inferences). Let's consider the Real numbers, it has axioms Let "a", "b" and "c" be real numbers These are the "minimums" we assume to be true in this system. These are "minimum" in the sense that everything else that is true about this number system can be derived from those axioms! Let's consider the following true identity which is not included in the axioms, but we can prove it using the axioms. We proceed: Before we proceed any further, you will have notice that the real numbers are not the only numbers that satisfies those axioms! For example the rational numbers also satisfy all the axioms. This leads to the abstract concept of a "field". In simple terms, a "field" is a number system that satisfies all those axiom. Let's define a "field" more carefully: A number system, "F", is a "field" if it supports + and × operations such that: Now, for M3, we do not let "b" be zero, since 1/0 has no meaning. However for the "M" axioms, we have excluded zero anyway. For interested students, the requirements of "closure", "identity", having "inverses" and "associativity" on an operation and a set are known as a . If "F" is a group with addition and "F"* is a group with multiplication, plus the distributivity requirement, "F" is a field. The above axioms merely state this fact in full. Note that the natural numbers are not a field, as M3 is generally not satisfied, i.e. not every natural number has an inverse that is also a natural number. Please note also that (-"a") denotes the additive inverse of "a", it doesn't say that (-a) = (-1)(a), although we can prove that they are equivalent. Example 1 Prove using only the axioms that 0 = -0, where -0 is the additive inverse of 0. Solution 1 Example 2 Let F be a field and "a" an element of F. Prove using nothing more than the axioms that 0"a" = 0 for all "a". Solution Example 3 Prove that (-"a") = (-1)"a". Solution 3 One wonders why we need to prove such obvious things (obvious since primary school). But the idea is not to prove that they are true, but to practise inferencing, how to logically join up arguments to prove a point. That is a vital skill in mathematics. Exercises. 1. Describe a field in which 1 = 0 2. Prove using only the axioms if u + v = u + w then v = w (subtracting u from both sides is not accepted as a solution) 3. Prove that if xy = 0 then either x = 0 or y = 0 4. In F-, the operation + is defined to be the difference of two numbers and the × operation is defined to be the ratio of two numbers. E.g. 1 + 2 = -1, 5 + 3 = 2 and 9×3 = 3, 5×2; = 2.5. Is F- a field? 5. Explain why Z6 (modular arithmetic modular 6) is not a field. Problem Set. 1. Prove for formula_82 2. Prove by induction that formula_83 3. Prove by induction where 4. Prove by induction formula_87 5. Prove that if x and y are integers and n an odd integer then formula_88 is an integer. 6. Prove that (n~m) = n!/((n-m)!m!) is an integer. Where n! = n(n-1)(n-2)...1. E.g. 3! = 3×2×1 = 6, and (5~3) = (5!/3!)/2! = 10. "Many questions in other chapters require you to prove things. Be sure to try the techniques discussed in this chapter." 

Conjugating 'to be'. In these cases, we use the correct form of "sein" for each situation. Please notice the final two sentences both use 'Sie', and we must look at the verb to determine the difference between 'she' and 'they'. In German, the English infinitive 'to be' is translated as "sein". This is the table of the forms of 'sein', with rough English translations. Note that in English, there are only three forms (am, is, are) while German has five (bin, bist, ist, sind, seid). Also, the verb conjugation of the two you-formals are always the exact same. German sein English to be Conjugating Normal Verbs. In these sentences, different verbs and endings are used. Note that the verb is always in second position. When conjugating normal verbs, use the endings shown below (a memory hook is the "best ten" endings). Note that in normal verbs, such as spielen and machen, ihr-form and er/sie/es-form are the same and the wir-form, sie (pl)-form and the formal are all the same as the infinitive. -en spielen - to play machen - to make/do Conjugating Irregular Verbs. In each of these sentences, we use an irregular verb. Irregularity occurs in the ich-form or the du-form and er/sie/es-forms. There are three types of irregularity. E in the first syllable. One form of irregularity occurs "sometimes" when the verb contains an 'e' in the first syllable. The change is simple: the du-form and er/sie/es forms both change the 'e' to an 'i.e.' or an 'i'. Two common examples are shown. Note that the er/sie/es-form and ihr-form are no longer the same. sehen - to see geben - to give Haben. A similar, yet different, change occurs in the verb "haben". As in the irregularity above, the du-form and er/sie/es-form change. haben - to have Verbs ending in Consonant-N. Some verbs change the ich-form for obvious reasons. "Wandern" and "basteln" are two examples. Both drop the first e in the ich-form. wandern - to hike basteln - to build Conjugating Modals. Modals are a new kind of verb. They are the equivalent to helping verbs in English. There are seven basic modals: können (can), mögen (like), dürfen (may), wollen (want), sollen (should), müssen (must), and möchten (would like). Möchten isn't technically a modal, but it acts like one in most aspects. Modals are conjugated very differently. The ich-form and er/sie/es-form are always alike and singular has a different verb in the first syllable (except in sollen and möchten). Below are the conjugations of the six basic modals and möchten. können - can mögen - like dürfen - may wollen - want sollen - should müssen - must möchten - would like Separable Verbs. Some verbs in German are separable: they have a prefix that can be separated from the base. When the verb is used with a modal, it regains the prefix at the end of the sentence. When it is the main verb of the sentence, the prefix is moved to the end of the sentence. An "example" in English would be the word "intake". When it is used as a verb, it becomes "take ... in". When it is used as an adjective or a noun, it becomes "intake" again. Two easy examples of separable verbs are "aussehen" and "mitkommen". Note that aussehen is also irregular. aussehen - to appear mitkommen - to come along/with 

=Force= A net force on a body causes a body to accelerate. The amount of that acceleration depends on the body's inertia (or its tendency to resist changes in motion), which is measured as its mass. When Isaac Newton formulated Newtonian mechanics, he discovered three fundamental laws of motion. Later, Albert Einstein proved that these laws are just a convenient approximation. These laws, however, greatly simplify calculations and are used when studying objects at velocities that are small compared with the speed of light. Friction. It is the force that opposes relative motion or tendency of relative motion between two surfaces in contact represented by f. When two surfaces move relative to each other or they have a tendency to move relative to each other, at the point (or surface) of contact, there appears a force which opposes this relative motion or tendency of relative motion between two surfaces in contact. It acts on both the surfaces in contact with equal magnitude and opposite directions (Newton's 3rd law). Friction force tries to stop relative motion between two surfaces in contact, if it is there, and when two surfaces in contact are at rest relative to each other, the friction force tries to maintain this relative rest. Friction force can assume the magnitude (below a certain maximum magnitude called limiting static friction) required to maintain relative rest between two surfaces in contact. Because of this friction force is called a self adjusting force. Earlier, it was believed that friction was caused due to the roughness of the two surfaces in contact with each other. However, modern theory stipulates that the cause of friction is the Coulombic force between the atoms present in the surface of the regions in contact with each other. Formula: Limiting Friction = (Friction Coefficient)(Normal reaction) Static Friction = the friction force that keeps an object at relative rest. Kinetic Friction = sliding friction Newton's First Law of Motion. This means, essentially, that acceleration does not occur without the presence of a force. The object tends to maintain its state of motion. If it is at rest, it remains at rest and if it is moving with a velocity then it keeps moving with the same velocity. This tendency of the object to maintain its state of motion is greater for larger mass. The "mass" is, therefore, a measure of the inertia of the object. In a state of equilibrium, where the object is at rest or proceeding at a constant velocity, the net force in every direction must be equal to 0. At a constant velocity (including zero velocity), the sum of forces is 0. If the sum of forces does not equal zero, the object will accelerate (change velocity over time). It is important to note, that this law is applicable only in non-accelerated coordinate systems. It is so, because the perception of force in accelerated systems are different. A body under balanced force system in one frame of reference, for example a person standing in an accelerating lift, is acted upon by a net force in the earth's frame of reference. Inertia is the tendency of an object to maintain its velocity i.e. to resist acceleration. Newton's Second Law of Motion. These two statements mean the same thing, and is represented in the following basic form (the system of measurement is chosen such that constant of proportionality is 1) : The product of mass and velocity i.e. "m"v is called the momentum. The net force on a particle is thus equal to rate change of momentum of the particle with time. Generally mass of the object under consideration is constant and thus can be taken out of the derivative. Force is equal to mass times acceleration. This version of Newton's Second Law of Motion assumes that the mass of the body does not change with time, and as such, does not represent a general mathematical form of the Law. Consequently, this equation cannot, for example, be applied to the motion of a rocket, which loses its mass (the lost mass is ejected at the rear of the rocket) with the passage of time. An example: If we want to find out the downward force of gravity on an object on Earth, we can use the following formula: Hence, if we replace m with whatever mass is appropriate, and multiply it by 9.806 65 m/s2, it will give the force in newtons that the earth's gravity has on the object in question (in other words, the body's weight). Newton's Third Law of Motion. This means that for every force applied on a body A by a body B, body B receives an equal force in the exact opposite direction. This is because forces can only be applied by a body on another body. It is important to note here that the pair of forces act on two different bodies, affecting their state of motion. This is to emphasize that pair of equal forces do not cancel out. There are no spontaneous forces. It is very important to note that the forces in a "Newton 3 pair", described above, can never act on the same body. One acts on A, the other on B. A common error is to imagine that the force of gravity on a stationary object and the "contact force" upwards of the table supporting the object are equal by Newton's third law. This is not true. They may be equal - but because of the second law (their sum must be zero because the object is not accelerating), not because of the third. The "Newton 3 pair" of the force of gravity (= earth's pull) on the object is the force of the object attracting the earth, pulling it upwards. The "Newton 3 pair" of the table pushing it up is that it, in its turn, pushes the table down. Equations. "To find Displacement" "To find Final Velocity" "To find Final Velocity" "To find Force when mass is changing" "To find Force when mass is a constant" 

Newtonian Gravity. Newtonian Gravity (simplified gravitation) is an "apparent" force (a.k.a. "pseudoforce") that simulates the attraction of one mass to another mass. Unlike the three fundamental (real) forces of electromagnetism and the strong and weak nuclear forces, gravity is purely attractive. As a force it is measured in newtons. The distance between two objects is measured between their centers of mass. Gravitational force is equal to the product of the universal gravitational constant and the masses of the two objects, divided by the square of the distance between their centers of mass. The value of the gravitational field which is equivalent to the acceleration due to gravity caused by an object at a point in space is equal to the first equation about gravitational force, with the effect of the second mass taken out. Gravitational potential energy of a body to infinity is equal to the universal gravitational constant times the mass of a body from which the gravitational field is being created times the mass of the body whose potential energy is being measured over the distance between the two centers of mass. Therefore, the difference in potential energy between two points is the difference of the potential energy from the position of the center of mass to infinity at both points. Near the earth's surface, this approximates: Potential energy due to gravity near the earth's surface is equal to the product of mass, acceleration due to gravity, and height (elevation) of the object. If the potential energy from the body's center of mass to infinity is known, however, it is possible to calculate the escape velocity, or the velocity necessary to escape the gravitational field of an object. This can be derived based on utilizing the law of conservation of energy and the equation to calculate kinetic energy as follows: Variables&lt;br&gt; Definition of terms A black hole is a geometrically defined region of space time exhibiting such large centripetal gravitational effects that nothing such as particles and electromagnetic radiation such as light may escape from inside of it. That is the escape velocity upon the event horizon is equivalent to the speed of light. General relativity is a metric theory of gravitation generalizing space time and Newton's law of universal gravitational attraction as a geometric property of space time. 

=Work = Work is equal to the scalar product of force and displacement. The scalar product of two vectors is defined as the product of their lengths with the cosine of the angle between them. Work is equal to force times displacement times the cosine of the angle between the directions of force and displacement. Work is equal to change in kinetic energy plus change in potential energy for example the potential energy due to gravity. Work is equal to average power times time. The Work done by a force taking something from point 1 to point 2 is Work is in fact just a transfer of energy. When we 'do work' on an object, we transfer some of our energy to it. This means that the work done on an object is its increase in energy. Actually, the kinetic energy and potential energy is measured by calculating the amount of work done on an object. The gravitational potential energy (there are many types of potential energies) is measured as 'mgh'. mg is the weight/force and h is the distance. The product is nothing but the work done. Even kinetic energy is a simple deduction from the laws of linear motion. Try substituting for v^2 in the formula for kinetic energy. Variables&lt;br&gt;  Definition of terms When work is applied to an object or a system it adds or removes kinetic energy to or from that object or system. More precisely, a net force in one direction, when applied to an object moving opposite or in the same direction as the force, kinetic energy will be added or removed to or from that object. Note that work and energy are measured in the same unit, the joule (J). Advanced work topics 

=Energy= Kinetic energy is simply the capacity to do work by virtue of motion. (Translational) kinetic energy is equal to one-half of mass times the square of velocity. (Rotational) kinetic energy is equal to one-half of moment of inertia times the square of angular velocity. Total kinetic energy is simply the sum of the translational and rotational kinetic energies. In most cases, these energies are separately dealt with. It is easy to remember the rotational kinetic energy if you think of the moment of inertia I as the "rotational mass". However, you should note that this substitution is not universal but rather a rule of thumb. Potential energy is simply the capacity to do work by virtue of position (or arrangement) relative to some zero-energy reference position (or arrangement). Potential energy due to gravity is equal to the product of mass, acceleration due to gravity, and height (elevation) of the object. Note that this is simply the vertical displacement multiplied by the weight of the object. The reference position is usually the level ground but the initial position like the rooftop or treetop can also be used. Potential energy due to spring deformation is equal to one-half the product of the spring constant times the square of the change in length of the spring. The reference point of spring deformation is normally when the spring is "relaxed," i.e. the net force exerted by the spring is zero. It will be easy to remember that the one-half factor is inserted to compensate for finite '"change in length" since one would want to think of the product of force and change in length formula_1 directly. Since the force actually varies with formula_2, it is instructive to need a "correction factor" during integration. Variables&lt;br&gt; Definition of terms 

=Momentum= Linear momentum. Momentum is equal to mass times velocity. Angular momentum. Angular momentum of an object revolving around an external axis formula_1 is equal to the cross-product of the position vector with respect to formula_1 and its linear momentum. Angular momentum of a rotating object is equal to the moment of inertia times angular velocity. Force and linear momentum, torque and angular momentum. Net force is equal to the change in linear momentum over the change in time. Net torque is equal to the change in angular momentum over the change in time. Conservation of momentum. Let us prove this law. We'll take two particles formula_3 . Their momentums are formula_4 . They are moving opposite to each other along the formula_5-axis and they collide. Now force is given by: According to Newton's third law, the forces on each particle are equal and opposite.So, Rearranging, This means that the sum of the momentums does not change with time. Therefore, the law is proved. Variables&lt;br&gt; Calculus-based Momentum. Force is equal to the derivative of linear momentum with respect to time. Torque is equal to the derivative of angular momentum with respect to time. 

Lektion Vier für Fortgeschrittene Vokabeln 4-3.  der Ärmelkanaltunnel Chunnel (England-France channel tunnel)  die Arbeit work  die Bibliothek library  die Buchhaltung accounting office  das Büro office  der Donnerstag Thursday  die Geschäftsbibliothek company (business) library  der Montag Monday  der Name name  der Schnellzug express train  das Sehen vision  die Versammlung meeting  das Wien Vienna (Austria)  das Wiedersehen reunion  die Woche week  das Zürich Zurich  alles klar all right, everything clear  am Montag on Monday  dann wenn at such time when  Darf ich... ? May I... ?  Es freut mich sehr It gives me pleasure  Guten Morgen! Good morning! ("greeting")  Ja, gewiss certainly, of course  vor Ende der Woche before the end of the week  Wiener Büro Vienna branch office  abhalten hold  abschließen complete  ankommen (kam an,  angekommen) arrive  fahren ride  geben give  kennen lernen meet, make acquaintance  müssen must ("aux.")  reisen travel  sehen see, look  tun do, accomplish  sich vorstellen introduce  werden will  würde would  bitte please  da there  durch through, by means of  endlich finally  gestern yesterday  nach to, towards  natürlich of course  mich myself ("reflexive")  mit with  schnell fast, quick, rapid  sofort directly, forthwith  wieder again, once again Grammatik 4-4 ~ Personal Pronouns: Accusative Case. Here are the personal pronouns in the accusative case:  * The accusative case is that of the object of a verb. Only transitive verbs take direct objects. The pronoun (and noun in two cases) object in each of these sentences is underlined in the German and the English:  "Können Sie mich verstehen?" Can you understand me?  "Ich kann Sie verstehen." I can understand you.  "Ich kann sie verstehen" I can understand (her or them).  "Ich kann ihn dir zurück kicken!" I can kick it back to you! Note the order of the pronouns in this last sentence. If the direct object (here: "ihn") is a personal pronoun, it precedes the dative ("dir"); if it were a noun, the dative would precede it, as in these sentences:  "Hier, ich kicke dir den Ball zu." Here, I kick the ball to you.  "Darf ich Ihnen meine Freundin vorstellen?" May I introduce my friend to you? Other uses of the accusative case in German will be explored in future lessons. Tables of the personal pronouns in all cases are summarized in Pronoun Tables. Grammatik 4-5 ~ Personal Pronouns in the Dative Case. Here are the personal pronouns in the dative case: The dative case is that of the indirect object of a verb. The pronoun indirect object of these sentences is underlined in the German and the English:  "Es geht mir gut" It goes (for) me well  "Wie geht es dir?" How goes it (for or with) you  "Und können Sie mir sagen...?" And can you tell me...?  "Karl gibt ihm den Ball" Karl gave him the ball.  "Wie geht es Ihnen?" How goes it (with) you? (How are you?) This last sentence is an example from Gespräch 1-2 using the polite form of 'you'. Whether singular or plural must be established by context. This next sentence translates with "ihnen" as 'them':  "Wie geht es ihnen?" How goes it with them? (How are they?) The meaning of "ihnen" (or "Ihnen") would have to come from context in a conversation. Another use of the dative case in German is after these prepositions: aus, bei, mit, nach, seit, von, zu. You will be introduced to the meanings of these prepositions over many future lessons rather than all at once, because some have many meanings in English. Indeed, because each language associates specific prepositions with many common sayings (and these often do not correspond in German and English), these "little" words can be troublesome for students. Nonetheless, you should memorize now the list of prepositions above to always remember their association with the dative case. Tables of the pronouns in all cases are summarized in Appendix 2. Word order in a German sentence with an indirect object depends upon whether that direct object is a pronoun or a noun. If the direct object is a noun, the dative precedes the accusative; if the direct object is a personal pronoun, the accusative precedes the dative: English sentence structure is similar. 

"This wikibook concerns about the history and development of thoughts, theology, and philosophy stemming from the Islamic faith throughout the centuries. For the political history of Islamic civilization, see History of Islamic Civilization." The history of Islam begins in the 7th century AD. 

Grammar - Past Participle (el participio). Spanish uses the past participle primarily for present perfect, past perfect, and other similar times. For -ar verbs form the past participle by adding -ado to the stem. For -er and -ir verbs add -ido: If the stem of an -er or -ir verb ends in one of the vowels -a, -e, or -o, the i of -ido gets an accent mark: There are a few verbs with an irregular past participle: As in English, the past participle can also be used as an adjective for a noun. In that case the ending has to match gender and number of the noun. Example: Finally, there are a few verbs with "both" a regular and an irregular past participle. In this case, the irregular past participle is used as an adjective, while the regular form is used for the verb tenses. Grammar - Present Perfect (el pretérito perfecto). The Spanish present perfect is formed by conjugating the auxiliary verb haber (= "to have") and adding the past participle of the verb. Here are a few examples of the Spanish Present perfect. Note that in Spanish the auxiliary verb haber and the past participle are never separated: Grammar - Pluperfect (el pretérito pluscuamperfecto). The Spanish pluperfect is formed by conjugating imperfect of haber (= "to have") and adding the past participle of the verb. Here are a few examples of the Spanish pluperfect. It is used to refer to an event that happened before another event in the past. As in the present perfect, the auxiliary verb haber and the past participle are never separated: 

 Lektion 2 "Fremde und Freunde" ~ Strangers and Friends Grammatik 2-1 ~ Introduction to Verbs. A is that part of speech that describes an action. Verbs come in an almost bewildering array of tenses, aspects, and types. For now, we will limit our discussion to verbs used in the present tense — i.e., describing an action occurring in the present. You should start to recognize that the form a verb takes is related to the subject of that verb: the verb form must match the person of the subject. This requirement is sometimes evident in English, but always so in German. Consider the following English and German sentences (the verb is "studieren" in every case): Several things are illustrated by these sentence pairs. First, all verbs in German follow the rule just stated that a verb form must agree with its subject. Starting in "Lektion 6" we will learn the verb forms associated with each person in German. Second, this rule in English applies mostly to the verb 'to be' (e.g., I am, you are, he is, etc.). In some English verbs, the 3rd person singular form is unique, often taking an 's' or 'es' ending: "I give at the office", but "He gives at the office" (and "She studies..." above). Finally, some German verbs are best translated with an English 'to be' verb form added. This is called the "progressive" form in English ('What are you studying?'), but it does not exist in German. Thus, a verb like "nennen" can best be translated as "to name" or "to call". The following example may make this clearer. In the present tense, the following statements in English: are all expressed in German in only one way: "Sie nennen die Firma, "Trans-Global"". And the question statement: 'Do they call the corporation, "Trans-Global"?' becomes, in German: "Nennen sie die Firma, "Trans-Global"?" Grammatik 2-2 ~ Pronouns in the Nominative Case. Most of the personal pronouns introduced in Lektion 1 are used as subjects of their verbs. These represent the nominative case in German (as in English). We will shortly learn three other cases in German: the accusative for direct objects, the dative for indirect objects, and the genitive for expressing possession. For now, remember that the singular personal pronouns in English (nominative case) are "I", "you", and "he/she/it" (1st, 2nd, and 3rd persons) and the nominative case is used as the subject of a verb. In German, these pronouns are rendered as "ich", "du", and "er/sie/es". In these example sentences, the subject of the verb is underlined: There are, of course, plural personal pronouns in the English nominative case: "we", "you", and "they"; and in German, these nominative case pronouns are "wir", "ihr", and "sie". These appear in the following examples (again, subject underlined): In both English and German, the 3rd person singular also has gender. As you will next learn, the 2nd person (person being addressed) in German has both familiar and polite (formal) forms. Further, it is worth repeating here — although introduced in "Grammatik 2-1" above and to be covered in detail in future lessons — that the verb form changes when the subject changes. That is, in German the verb form must match the subject of a sentence. Here are some examples; compare with the previous three example sentences above and note how the verb form changed to match the sentence subject (subject and verb underlined): In the last example, the English verb form ('have') also changed based upon the subject of the sentence. Gespräch 2-1 ~ "Die Geschäftsleute". In this conversation, although the subject matter is basically casual, a more formal form of German is being used intoning respect between coworkers in an office setting. The polite form is expressed by the pronouns as explained below (Grammatik 2-3). Vokabeln 2-1.  die Anleitungen instructions  das Deutsch German (language) (more common is "die deutsche Sprache")  der Fremde foreigner, stranger  die Firma company, firm, business concern  die Frage question  die Geschäftsleute business people ("die Leute" = people)  der Hauptsitz head office ("das Haupt" = head or chief)  der Tag day, daytime  aus England from England  Das ist richtig! That is right!  Frau Baumann Ms. Baumann  Herr Schmidt Mr. Schmidt  zu Besuch visiting  arbeiten work  getroffen (have) met (past participle of "treffen")  nennen name, call  alle all  an at  Ihnen (with "or" to) you (polite form)  heute today  ihr you (plural), you all  ja yes  nein no  richtig correct  sie they (note: also "she")  Sie you (polite form)  wir we Grammatik 2-3 ~ Familiar and Polite Pronoun Forms. Many pronouns were introduced in Lesson 1. In "Grammatik 2-1" and "Gespräch 2-1" we have been presented with the following additional pronouns:  "Ihnen" – (to) you (2nd person singular, dative case)  "ihr" – you (2nd person, plural, nominative case)  "sie" – they (3rd person, plural, nominative case)  "Sie" – you (2nd person, singular, nominative case)  "wir" – we (1st person, plural, nominative case) In the conversations between friends presented in "Gespräche 1-1" and "1-2" (Lektion 1) the familiar form of the personal pronouns (e.g., "du", "dir") was used. However, German also has a polite or formal form of some of these personal pronouns. The polite form is used in conversations between strangers and more formal situations, as illustrated in the "Gespräch 2-1": greetings between business associates. The polite form is always first-letter capitalized in German, which can be helpful in differentiating "Sie" (you) from "sie" (she and they); "Ihnen" (you) from "ihnen" (them). However, you will soon learn that the form of the verb (see "Grammatik 2-3" below) is most telling, as shown by these example pairs using the verb, "haben" (have): Because the first letter in a sentence is always capitalized, we cannot determine (without the verb form) whether the second and third examples begin with "sie" ('she' or 'they') or with "Sie" (polite 'you'); a problem that would also exist in conversation. The fourth example, where subject and verb are reversed in a question, demonstrates the pronoun 'they'; compare it with the polite 'you' in the first example. It is relatively easy for an English speaker to appreciate how context, especially in conversation, overcomes confusion considering that English has fewer forms for these pronouns than German. However, this fact does present some difficulty when learning German, since improper use of a pronoun may just create confusion in speaking or writing German. Gespräch 2-2 ~ "Die Geschäftsmänner". &lt;br&gt; &lt;br&gt; &lt;br&gt; Vokabeln 2-2.  die Bundesrepublik Deutschland Federal Republic of Germany  die Geschäftsmänner businessmen ("die Geschäftsleute" is preferred)  Großbritannien Great Britain (technically "Vereinigtes Königreich"  "von Großbritannien und Nordirland")  der Morgen morning  die Übersetzung translation  bis morgen until tomorrow  Guten Morgen! Good morning (greeting)  nicht so gut not so well  so viel so much  Wie bitte? How is that?  zu viel too much  bis until  kein no (in the sense on "none")  müde tired  nicht not  sich each other  warum ? why ? Grammatik 2-4 ~ Personal pronoun gender. In both English and German the 3rd person personal pronouns have gender (Grammatik 1-3). However, in English, the pronoun "it" is used for most inanimate or non-living things. There are a few exceptions: a ship might be referred to as "she". However, in German, the 3rd person personal pronoun reflects the gender of the noun (antecedent) referred to by the pronoun. For examples: The following table summarizes these gender relationships: Übersetzung 2-1. You may, at this point, try the flash cards developed for Level I German. This set has a few words and concepts not yet presented in Level II, but for the most part can be very helpful in enhancing your vocabulary. Go to FlashcardExchange.com. Translate the following sentences into German. Pay attention to whether familiar or polite form of the pronoun is requested: 

"The Once and Future King" - Minor Characters See also Major Characters. As many characters have only a given name, or take "of Somewhere" as their last name, the characters are listed alphabetically by first name. Note that characters in "The Once and Future King" may not exactly match with traditional Arthurian legend. Sir Ector - Owns the Castle of the Forest Sauvage. He adopts and cares for the Wart . He is the father of Kay. Sir Kay - Son of Sir Ector and friend of the Wart. He tries to act superior over the Wart. He follows Arthur on one adventure: the rescue mission at Morgan le Fay's castle (see "../The Sword in the Stone/"). When Arthur brings Kay the Sword in the Stone, Kay falsely claims to have drawn the sword himself, but soon gives the credit to Arthur. He later becomes a knight of the Round Table. Morgan le Fay - The most evil of the Cornwalls, and the strongest in magic. "Le Fay" means "the Fairy". She is the antagonist in the major adventure in "../The Sword in the Stone/". Her sisters are Elaine (not mentioned) and Morgause. 

Lektion Drei für Fortgeschrittene Gespräch 3-3 ~ "Mach dir keine Sorgen!". Beim Ballspielen macht Karl sich Sorgen um die Uhrzeit. Vokabeln 3-3.  das Ballspiel ball game  die Minute minute  das Motorrad motorcycle  die Sorge, die Sorgen problem(s), worry(-ies)  das Viertel quarter, one-fourth  die Woche week  die Wohnung apartment  mach dir keine Sorgen! do not worry!  nach Hause gehen go home  kicken kick  zurückkicken kick back, return kick  beim when, while (usually, "at the")  danach after that  dein your  erst only  halb half  jetzt now  komisch comical, funny  mein my  schon already  zurück back  warum why (interrogative) Grammatik 3-5 ~ Numbers. Gender of Ordinals. Ordinal numbers are adjectives, and therefore have forms for each of the three genders in German. The forms are derived from the feminine form (as introduced in the beginning of Lesson 3) by adding an 'r' (masculine) or an 's' (neuter). Thus: "erste" (feminine), "erster (masculine), and "erstes (neuter). Examples: Grammatik 3-6 ~ Expressions of Time. Idioms used in Telling Time. As in English, there are a number of idiomatic phrases associated with giving or telling time. For example, note that the half hour is given as approaching the next hour. The German preposition, "um", is used to mean "at" a given time. Periods of the Day. There are a number of adverbial phrases used in German to denote time periods during the day. Common ones are listed here: Additional Notes. The first sentence in Gespräch 3-3 uses "Beim Ballspielen" in the sense of "during the ball game" or "while playing ball". "Beim" is a contraction of "bei dem" or "at the". However, "das Ballspiel" is a noun that represents an action ("playing with a ball"), so it is correct to use "beim" in the sense intended here. It is not the most beautiful way of saying this—but is correct. With the infinitive of a verb you can use "beim" too: "Beim Spielen" means "while playing". This form is more common in modern German language. Vokabeln 3-4.  der Abend evening  der Himmel heaven  der Mittag noon, noontime  der Morgen, die Morgen morning(s)  der Nachmittag afternoon  die Nacht night  der Tag, die Tage day(s)  abreisen depart (from a trip)  auf for (duration), after  gegen towards, about, approximately  letzt(er) last  ungefähr (at) about, approximately Note that "morgen" does not change in plural; thus, "Die Morgen" = "the mornings". It is uncommon to use it in plural, unless as a measure of land "Vier Morgen Land" = "four 'morgens of land". For a plural use of "mornings", it is better to substitute "die Vormittage". Andere Wörter 3A. Using these additional vocabulary words, you may be able to restate Gespräch 3-3 above, altering the meaning (or time of day) of the conversation.  die Hälfte half  die Viertelstunde quarter of an hour Übersetzung 3-2. Translate the following sentences into German: 

Oxymercuration is a process by which water is added to an alkene through treatment of the alkene with mercury(II) acetate [Hg(O2CCH3)2, usually abbreviated Hg(OAc)2] in an aqueous tetrahydrofuran (THF) solvent. Rarely this reaction is referred to as Oxymercuration/reduction. This reaction looks just like Markovnikov. addition of water across a double bond, however there is no carbocation intermediate, so there is no rearrangement. You will learn that NaBH4 is a common reducing agent. 

They are powers of 3: 3 9 27 81 243 729 2,187 ... 

The Text Formatting elements give logical structure to phrases in your HTML document. This structure is normally presented to the user by changing the appearance of the text. We have seen in the Introduction to this book how we can "emphasize" text by using codice_1 tags. Graphical browsers normally present emphasized text in italics. Some Screen readers, utilities which read the page to the user, may speak emphasized words with a different inflection. A common mistake is to tag an element to get a certain "appearance" instead of tagging its "meaning". This issue becomes clearer when testing in multiple browsers, especially with graphical and text-only browsers as well as screen readers. You can change the default presentation for any element using Cascading Style Sheets. For example, if you wanted all emphasized text to appear in red normal text you would use the following CSS rule: In this section, we will explore a few basic ways in which you can markup the logical structure of your document. Emphasis. HTML has elements for two degrees of emphasis: An example of emphasized text: It is essential not only to guess but actually &lt;em&gt;observe&lt;/em&gt; the results. An example rendering: An example of strongly emphasized text: Let us now focus on &lt;strong&gt;structural markup&lt;/strong&gt;. An example rendering: Preformatted text. Preformatted text is rendered using fixed-width font, and without condensing multiple spaces into one, which results in preserved spacing. Newlines are rendered as newlines, unlike outside preformatted text. HTML markup in the preformatted text is still interpreted by browsers though, meaning that "a" will still be rendered as "a". To create preformatted text, start it with &lt;pre&gt; and end it with &lt;/pre&gt;. An example: &lt;pre&gt; &lt;/pre&gt; The resulting rendering: Omitting the preformatting tags will cause the same text to appear all in one line: Special Characters. To insert non-standard characters or characters that hold special meaning in HTML, a character reference is required. For example, to input the ampersand, "&amp;", "&amp;amp;" must be typed. Characters can also be inserted by their ASCII or Unicode number code. Abbreviations. Another useful element is codice_4. This can be used to provide a definition for an abbreviation, e.g.  &lt;abbr title="HyperText Markup Language"&gt;HTML&lt;/abbr&gt; Graphical browsers often show abbreviations with a dotted underline. The codice_5 appears as a tooltip. Screen readers may read the codice_5 at the user's request. Note: very old browsers (Internet Explorer version 6 and lower) do not support codice_4. Because they support the related element codice_8, that element has been commonly used for all abbreviations. An acronym is a special abbreviation in which letters from several words are pronounced to form a new word (e.g. radar - Radio Detection And Ranging). The letters in HTML are pronounced separately, technically making it a different sort of abbreviation known as an initialism. Discouraged Formatting. HTML supports various formatting elements whose use is discouraged in favor of the use of cascading style sheets (CSS). Here's a short overview of the discouraged formatting, so that you know what it is when you see it in some web page, and know how to replace it with CSS formatting. Some of the discouraged elements are merely discouraged, others are deprecated in addition. Cascading Style Sheets. The use of style elements such as &lt;b&gt; for bold or &lt;i&gt; for "italic" is straight-forward, but it couples the presentation layer with the content layer. By using Cascading Style Sheets, the HTML author can decouple these two distinctly different parts so that a properly marked-up document may be rendered in various ways while the document itself remains unchanged. For example, if the publisher would like to change cited references in a document to appear as bold text as they were previously "italic", they simply need to update the style sheet and not go through each document changing &lt;b&gt; to &lt;i&gt; and vice-versa. Cascading Style Sheets also allow the reader to make these choices, overriding those of the publisher. Continuing with the above example, let's say that the publisher has correctly marked up all their documents by surround references to cited material (such as the name of a book) in the documents with the &lt;cite&gt; tag: &lt;cite&gt;The Great Gatsby&lt;/cite&gt; Then to make all cited references bold, one would put something like the following in the style sheet: Later someone tells you that references really need to be italic. Before CSS, you would have to hunt through all your documents, changing the &lt;b&gt; and &lt;/b&gt; to &lt;i&gt; and &lt;/i&gt; (but being careful *not* to change words that are in bold that are not cited references). But with CSS, it's as simple as changing one line in the style sheet to 

GCSE Science/Electricity The motor effect is the term used when a current-carrying wire in the presence of a magnetic field experiences a force. A simple experimental demonstration will show you that this is true. Place a wire that is connected to a power pack in between the poles of a horseshoe magnet. Turn on the power and the wire moves. Often the movement is only very slight because a typical horseshoe magnet is not very strong. The force depends on a number of things: The force obeys the formula: formula_1 The first two points are pretty much obvious, so let's look at the third point in a little more detail. The magnetic field of a horseshoe magnet points pretty much in a straight line from the north pole to the south pole. If the wire cuts this field at right angles the resulting force will be a maximum. If the wire runs parallel to the field, from the north the south pole or vice versa, the wire will still experience the motor effect. However, the net result of this force is along the wire, not perpendicular to it, thus the wire does not turn. Having said that, at GCSE you only really need to think about situations where the field and the wire cut each other at right angles. The force is always at right angles to both the field and the current flowing in the wire. This means that if you draw the direction of the magnetic field and the wire on a piece of paper the force will be out of the plane of the paper pointing straight up or down {More about how you work out which way later). Look at the diagram above. For simplicity, only the two ends of the horseshoe magnet have been drawn. Also the power pack and connecting wires are not shown either. The magnetic field is going into the screen. The current is traveling from right to left. The black line represents the force, and therefore the direction that the wire moves. Fleming's Left Hand Rule. This rule allows you to work out which direction the force will point in. Arrange your left hand with your thumb, first finger, and second fingers all pointing at right angles to one another. Remember you have to use your left hand for this! Although you may use your right hand so long as you swap the direction of current with the force Q4) A student uses her right hand instead of her left. What effect will that have on the force she works out? You will recall that a current carrying wire is surrounded by its own magnetic field. The diagram below shows the wire end on in a magnetic field of two magnets that are NS facing. The field due to the magnets is shown in blue, and the field due to the current in the wire is shown in black. Notice the direction of the two fields as shown by the arrows. On top of the wire the fields are both going in the same direction. They add up making an overall strong field. Underneath the wire, they go in opposite directions. They cancel each other out to some extent making an overall weaker field. The new field is shown in the diagram below. See how the field above the wire is stronger. The lines are closer together. Below the wire the field is weaker (due to partial canceling out) the field lines are further apart. The force pushes the wire downwards, away from the strong field into the weak field. It's as if the field lines try to repel each other. They don't like being squashed together and try to straighten out. They also act as if they are made of elastic bands, they don't like being stretched out of shape. (This is just a model of what's going on. The lines aren't real, they don't actually try to push each other away, but I find it a way of helping me understand what's going on. If it doesn't help you, don't use it) A simple electric motor. Ok, so far we have been looking at the force that results when we put a current carrying wire in a magnetic field. In this section we will look at a practical use for this force.As you have probably already guessed from the name of this page, the practical use is going to be an electric motor. Look at the diagram above. A rectangular loop of wire is sitting inside a magnetic field. We can consider the current in the four sections of the loop and work out which way the force acts. The net result of these different forces is that there will be a turning moment that makes the coil rotate by 90°. At that point the upwards and downwards forces will be acting along the same line and the coil will stop turning. Another way to think about it is to consider the loop as a tiny little one turn solenoid. The solenoid will have a little north pole and a little south pole and will therefore move until its north pole lines up with the south pole of the magnet on the right, and its south pole lines up with the north pole of magnet on the left. This is all very interesting but not much use as a motor. We want something that keeps turning all the time the current flows. They way this is achieved is by the use of the commutator – a circular metal ring that is split into two halves. The ends of the wire loop turn around inside the commutator. They are in electrical contact with it. One side of the commutator is connected to the positive output of a power pack or battery . the other half of the commutator is connected to the negative. Let's look at what happens as the coil turns inside the commutator: Q7)A student sets up an electric motor and turns it on. The coil turns clockwise. List two ways she could reverse the direction. «Uses of electromagnets | Induction» 

In the chapter Variables we looked at the primitive data types. However "advanced" data types allow us greater flexibility in managing data in our program. Structs. Structs are data types made of variables of other data types (possibly including other structs). They are used to group pieces of information into meaningful units, and also permit some constructs not possible otherwise. The variables declared in a struct are called "members". One defines a struct using the codice_1 keyword. For example: struct mystruct {  int int_member;  double double_member;  char string_member[25]; } struct_var; codice_2 is a variable of type codice_3, which we declared along with the definition of the new codice_3 data type. More commonly, struct variables are declared after the definition of the struct, using the form: struct mystruct struct_var; It is often common practice to make a "type synonym" so we don't have to type "struct mystruct" all the time. C allows us the possibility to do so using a codice_5 statement, which aliases a type: typedef struct { } Mystruct; The codice_1 itself is an "incomplete" type (by the absence of a name on the first line), but it is aliased as codice_7. Then the following may be used: Mystruct struct_var; The members of a struct variable may be accessed using the member access operator codice_8 (a dot) or the indirect member access operator codice_9 (an arrow) if the struct variable is a pointer: struct_var.int_member = 0; struct_var-&gt;int_number = 0; // this statement is equivalent to: (*struct_var).int_number = 0; Structs may contain not only their own variables but may also contain variables pointing to other structs. This allows a recursive definition, which is very powerful when used with pointers: struct restaurant_order {  char description[100];  double price;  struct restaurant_order *next_order; This is an implementation of the linked list data structure. Each node (a restaurant order) is pointing to one other node. The linked list is terminated on the last node (in our example, this would be the last order) whose codice_10 variable would be assigned to codice_11. A recursive struct definition can be tricky when used with codice_5. It is not possible to declare a struct variable inside its own type by using its aliased definition, since the aliased definition by codice_5 does not exist before the codice_5 statement is evaluated: typedef struct Mystruct {  struct Mystruct *pointer; // Mystruct *pointer; would cause a compile-time error } Mystruct; The size of a struct type is at least the sum of the sizes of all its members. But a compiler is free to insert padding bytes between the struct members to align the members to certain constraints. For example, a struct containing of a char and a float will occupy 8 bytes on many 32bit architectures. Unions. The definition of a union is similar to that of a struct. The difference between the two is that in a struct, the members occupy different areas of memory, but in a union, the members occupy the same area of memory. Thus, in the following type, for example: union {  int i;  double d; } u; The programmer can access either codice_15 or codice_16, but not both at the same time. Since codice_15 and codice_16 occupy the same area of memory, modifying one modifies the value of the other, sometimes in unpredictable ways. This is also the main reason that unions are rarely seen in practice. The size of a union is the size of its largest member. Enumerations. Enumerations are artificial data types representing associations between labels and integers. Unlike structs or unions, they are not composed of other data types. An example declaration: enum color {  red,  orange,  yellow,  green,  cyan,  blue,  purple, } crayon_color; In the example above, red equals 0, orange equals 1, ... and so on. It is possible to assign values to labels within the integer range, but they must be a literal. Similar declaration syntax that applies for structs and unions also applies for enums. Also, one "normally" doesn't need to be concerned with the integers that labels represent: enum weather weather_outside = rain; This peculiar property makes enums especially convenient in switch-case statements: enum weather {  sunny,  windy,  cloudy,  rain, } weather_outside; switch (weather_outside) { case sunny:  wear_sunglasses();  break; case windy:  wear_windbreaker();  break; case cloudy:  get_umbrella();  break; case rain:  get_umbrella();  wear_raincoat();  break; Enums are a simplified way to emulate associative arrays in C. 

=Review of logs= Been a while since you used logs? Here is a quick refresher for you. The log (short for logarithm) of a number N is the exponent used to raise a certain "base" number B to get N. In short, formula_1 means that formula_2. Typically, logs use base 10. An increase of "1" in a base 10 log is equivalent to an increase by a power of 10 in normal notation. In logs, "3" is 100 times the size of "1". If the log is written without an explicit base, 10 is (usually) implied. Another common base for logs is the trancendental number formula_3, which is approximately 2.7182818... Since formula_4, these can be more convenient than formula_5. Often, the notation formula_6 is used instead of formula_7. The following properties of logs are true regardless of whether the base is 10, formula_3, or some other number. &lt;br&gt; Adding the log of A to the log of B will give the same result as taking the log of the product A times B. &lt;br&gt;&lt;br&gt; Subtracting the log of B from the log of A will give the same result as taking the log of the quotient A divided by B. &lt;br&gt;&lt;br&gt; The log of (A to the Bth power) is equal to the product (B times the log of A). &lt;br&gt;&lt;br&gt; A few examples: &lt;br&gt; log(2) + log(3) = log(6) &lt;br&gt; log(30) – log(2) = log(15) &lt;br&gt; log(8) = log(23) = 3log(2) 

=Electricity= The force resulting from two nearby charges is equal to k times charge one times charge two divided by the square of the distance between the charges. This what force of attraction between to charged particle says, according to coulomb's law. The electric field created by a charge is equal to the force generated divided by the charge. Electric field is equal to a constant, “k”, times the charge divided by the square of the distance between the charge and the point in question. Electric potential energy is equal to a constant, “k” multiplied by the two charges and divided by the distance between the charges. Variables&lt;br&gt; Electricity acts as if all matter were divided into four categories: Charges are positive (+) or negative (-). Any two like charges repel each other, and opposite charges attract each other. Electric fields. A charge in an electrical field feels a force. The charge is not a vector, but force is a vector, and so is the electric field. If a charge is positive, then force and the electric field point in the same direction. If the charge is negative, then the electric field and force vectors point in opposite directions. A point charge in space causes an electric field. The field is stronger closer to the point and weaker farther away. Electricity is made of subatomic particles called Electrons and so are Electric Fields and Magnetic Fields. One must also note that electrical fields come under the category of spherical fields as the inverse square law may be applied to the electrical field. This means that the electrical force, exhibited by the electrical field emitted by the subatomic electron charge (-), acting upon a body is inversely proportional to the distance between the center point of the electric field (subatomic electron) and the body on which the electric force is acting upon. An electric circuit is composed of conducting wires (through which an electric current flows through), a key or switch which is utilized to open and close the circuit, components which transfer electrical energy to a form of energy required by the component and an electromotive source (such as a voltaic cell). A voltaic cell is an electromotive source in which are present two plates, zinc and copper, placed in dilute sulphuric acid. Whence the circuit is closed the zinc reacts with the sulphuric acid to produce zinc sulphate. The electromotive force which discharges the electrical energy in the electric current is considered to be originated on the surface of the zinc plate in the voltaic cell. However, depending upon the cell, closing the circuit gives rise to polarization, accumulation of hydrogen bubbles on the surface of the copper plate which seriously interferes with the movement of electricity and reduces the magnitude of the electromotive force. For this reason Leclanché cells are utilized. Consisting of similar characteristics as that of the voltaic cell however a Mage difference is present. Instead of the use of copper plates, a carbon plate is used. For this reason, magneze dioxide may be placed on the carbon to react to form a compound which whence in contact with hydrogen bubbles will turn the hydrogen into water, hence increasing the size of the electromotive force produced by the cell. The resistance encountered in conducting wires: Inversely proportional to the diameter of the conducting wire. Directly proportional to the length of the conducting wire. Varies with different substances. Varies with temperature of the conducting wire. In order to maintain a constant flow of an electric current a constant expenditure of chemical or mechanical energy is required. An electric current is accompanied by an electric field and a magnetic field. A device employed into determining the presence of an electric current is known as a galvanoscope. The conducting wire through which the electric current flows through is he led over and parallel to the galvanoscope the magnetoscope preset inside of the galvanoscope being deflected in the opposite direction to which the electric current flows in. So with the aid of a galvanoscope one may not only deduce the magnetic properties of an electric current, the exhibition of a magnetic field, but the direction in which the current flows through. An electromotive force may also be generators by a dynamo. A rotating magnet present inside of a helix. The magnetic properties of electric currents may be used to construct magnets. An electomagnet is commonly described as a mass of iron on which is placed a helix/solenoid through which flows an electric current. The magnetic field emitted by the electric current is increased if the solenoid is placed around a magnetic mass of iron or any other substance possessing magnetic properties, that is the magnetic field of the iron is added to that of the electic current producing a more powerful magnetic field. Conductors may be arranged in two variants. Series and parallel circuits. In series, the current passes through each conductor in turn, where Ohm's law changes to I = nE/(nr + R), where I is the current intensity, n is the number of cells arranged in series in the circuit, E is the electromotive force applied to the circuit, r is the internal resistance ( the resistance the current that is produced in the cell experiences whence passing from the zinc plate to the copper or carbon plate through the sulphuric acid ) and R is the external resistance. 

cmavo.  doi DOI vocative marker  generic vocative marker; identifies intended listener;  elidable after COI  coi COI greetings  vocative: greetings/hello  be'e COI request to send  vocative: request to send/speak  je'e COI roger  vocative: roger (ack) - negative acknowledge; used  to acknowledge offers and thanks  je'enai COI* negative acknowledge  vocative: roger (ack) - negative acknowledge;  I didn't hear you  co'o COI partings  vocative: partings/good-bye  re'i COI ready to receive  vocative: ready to receive - not ready to receive  re'inai COI* not ready to receive  vocative: ready to receive - not ready to receive  zo ZO 1-word quote  quote next word only; quotes a single Lojban word  (not a cmavo compound or tanru)  ma KOhA7 sumti ?  pro-sumti: sumti question (what/who/how/why/etc.);  appropriately fill in sumti blank  la LA that named  name descriptor: the one(s) called ... ; takes name  or selbri description  mi KOhA3 me  pro-sumti: me/we the speaker(s)/author(s); identified  by self-vocative doi. In English poetry sometimes people say O before names, but usually people just put the name at the beginning or the end of a sentence separated from the rest by a comma or as a sentence of its own. In English, people also say "Hey" before names. This carries the added effect of getting someone's attention. If I felt like saying something completely useless, I would say {doi do}, identifying the listener as the person listening. "Can someone say {doi} by itself (without a name following like I can with the rest of the vocatives)?" coi. Near the beginning of a conversation or upon noticing another Lojban speaker to talk to, a Lojban speaker will often say {coi}. The listener usually replies {coi}. This has about the same meaning as the English words "Hi", "Hello", "Howdy", and "Greetings". As with most of the vocatives, I can follow it with a name. It must either have la or a mandatory pause between {coi} and the name. A common Lojban phrase {coi ro do} means "Hi, all.". be'e. A Lojban speaker not speaking, wanting to speak, and not sure if other speakers will allow him/her to speak will sometimes say {be'e} to ask permission. In English, people usually only ask to speak on formal occasions like the courtroom and in large meetings. In school students raise their hands for a turn to speak. When people say "Excuse me.", they don't expect someone to deny them from speaking. They also ask "May I have a word with you?" Like most vocatives, I can follow {be'e} with the listener. While I see {be'enai} as correct grammar, I can't figure out a meaning for it beyond "Ha ha, I know a cmavo you don't!" "To me, {be'enai} seems fairly obvious: "Permission not to speak?" which could be used in a fairly similar way to "No comment!" but asking permission not to comment rather than refusing." je'e. In English, people acknowledge "Thank you.", with "Welcome.". They say "I heard." to acknowledge they heard the last utterance. They say "I understand." To point out they understand the last utterance. "Okay." also acknowledges the last utterance, but carries more hint of agreement than {je'e}. To show that nearly anyone would understand an utterance, English speakers sometimes use the word "Duh!". In English people say "Huh?" or "Wha?" or "What?" to denote similar meanings. co'o. A Lojban speaker says {co'o} to another Lojban speaker shortly before ending conversation or close proximity to the second Lojban speaker. The second Lojban speaker often replies {co'o}. As with most other vocatives, I can follow {co'o} with a name. English phrases with similar meanings include "Goodbye.", "So long.", and "See you later.". re'i. This word signifies that the speaker can listen to the listener. I can follow {re'i} with a name like most vocatives to signify the listener. English phrases with similar meanings include "Go ahead.", "I'm all ears.", and "I'm listening." la daniel. cusku zo be'e .i la kleir. cusku zo re'i Daniel says, "May I have a word with you?". Claire says, "Go ahead.". In English {re'inai} translate to "Sorry, I'm busy.", "I'm not accepting visitors.", and "I can't listen right now." zo. In English, I might put a word in quotes italics to quote it. If I only wanted to quote one word, I might say "the word foo". I can talk about the word {zo} as a word by saying {zo zo}. {zo nai} quotes the word {nai} instead of meaning the opposite of {zo}. ma. In asking another Lojban speaker for information, a Lojban speaker might say give a bridi describing their information with a blank for the listener to insert their answer. The asker would signify that blank with {ma}. The answerer would probably not repeat the entire le bridi, but instead just the information to replace {ma} with. In English, a person might say "What" or "blank" instead. la. zo la gadri mi. If I wanted to say something that didn't mean much. I would say, {mi'e mi}, pointing out the speaker as myself. "Does {minai} mean {do}?" "I believe that {minai} would mean anyone but yourself." gismu.  smuni mun smu meaning  x1 is a meaning/interpretation of x2  recognized/seen/accepted by x3  [referential meaning (=selsni, snismu)]; (cf. jimpe,  sinxa, valsi, tanru, gismu, lujvo, cmavo, jufra)  gismu gim gi'u root word  x1 is a (Lojban) root word expressing relation x2  among argument roles x3, with affix(es) x4  [gismu list, if physical object (= (loi) gimste);  referring to the mental construct (e.g. propose adding  a new gismu to the gismu list = gimpoi, gimselcmi,  gimselste)]; (cf. cmavo, cmene, lujvo, smuni, sumti,  tanru, valsi)  lujvo luv jvo affix compound  x1 (text) is a compound predicate word with meaning  x2 and arguments x3 built from metaphor x4  (cf. stura, cmavo, gismu, rafsi, smuni)  tanru tau phrase compound  x1 is a binary metaphor formed with x2 modifying x3,  giving meaning x4 in usage/instance x5  (x2 and x3 are both text or both si'o concept)  (cf. gismu, smuni)  sumti sum su'i argument  x1 is a/the argument of predicate/function x2 filling  place x3 (kind/number)  (x1 and x2 are text); (cf. bridi, darlu, gismu)  bridi bri predicate  x1 (text) is a predicate relationship with relation  x2 among arguments (sequence/set) x3  [also: x3 are related by relation x2 (= terbri  for reordered places)]; (x3 is a set completely  specified); (cf. sumti, fancu)  cmavo ma'o structure word  x1 is a structure word of grammatical class x2,  with meaning/function x3 in usage (language) x4  [x4 may be a specific usage (with an embedded language  place) or a massified language description; x3 and x4  may be merely an example of cmavo usage or refer to  an actual expression; cmavo list, if physical object  (= (loi) ma'oste); referring to the mental construct  (e.g. propose adding a new cmavo to the cmavo list =  ma'orpoi, ma'orselcmi, ma'orselste)]; (cf. gismu,  lujvo, gerna, smuni, valsi)  cmene cme me'e name  x1 (quoted word(s)) is a/the name/title/tag of x2  to/used-by namer/name-user x3 (person)  [also: x2 is called x1 by x3 (= selcme for reordered  places)]; (cf. cmavo list me'e, gismu, tcita,  valsi, judri) Exercises. In Exercises 1 to 11, translate the Lojban sentence to English.&lt;br&gt; 1. coi. braun. e'apei mi ba te cmene do lu la bab. li'u&lt;br&gt; 2. be'e meiris. mi'o ca cliva&lt;br&gt; 3. ko'a cusku lu le lujvo cu lujvo li'u .i ko'e cusku zo je'e&lt;br&gt; 4. ko'a cusku lu mi cmavo li'u .i ko'e cusku lu je'enai li'u&lt;br&gt; 5. co'o .ednas. ko ba klama mi&lt;br&gt; 6. la daniel. cusku zo be'e .i la kleir. cusku zo re'i&lt;br&gt; 7. mi cusku lu be'e. reitcel. li'u le nanmu pe la reitcel. cusku lu la reitcel. cusku lu re'inai li'u li'u&lt;br&gt; 8. mi cusku zo bridi .i mi na cusku tanru&lt;br&gt; 9. mi cusku lu lu smuni cmene li'u tanru ma ma li'u .i la keven. cusku lu smuni cmene li'u tanru zo smuni zo cmene li'u&lt;br&gt; 10. mi'e la mark.&lt;br&gt; 11. le prenu pe mi te cmene mi lu la brai'an. li'u In Exercises, 100 to 131, translate le jufra (the sentence) given.&lt;br&gt; 100. le se bridi cu smuni le selbri mi .i&lt;br&gt; 101. zo coi smuni zoi gy. hello .gy mi .i&lt;br&gt; 102. gy. goodbye .gy smuni zo co'o mi .i&lt;br&gt; 103. lu mi'e. brai'an. li'u smuni lu mi me la brai'an. li'u mi .i&lt;br&gt; 104. fi la'o smuni. George Bush .smuni goi ko'a smuni fa ko'a la'o smuni. Mr. President .smuni .i&lt;br&gt; 105. le ve gismu cu smuni le rafsi mi .i&lt;br&gt; 106. zoi ly. gi'u .ly ve gismu zo gismu .i&lt;br&gt; 107. zo smuni gismu .i&lt;br&gt; 108. zo gismu gismu fi le te gismu .i&lt;br&gt; 109. zo selbri gismu le du'u se bridi .i&lt;br&gt; 110. zo selbri lujvo lu se bridi li'u .i&lt;br&gt; 111. zo selbri lujvo zo'e le te lujvo be fa zo bridi .i&lt;br&gt; 112. zo selbri se ve se lujvo lu se bridi li'u .i&lt;br&gt; 113. zo selma'o lujvo fo lu se cmavo li'u .i&lt;br&gt; 114. zo selma'o lujvo fi le se cmavo .i&lt;br&gt; 115. zo selma'o lujvo lu se cmavo li'u le ve cmavo lu se cmavo li'u .i&lt;br&gt; 116. zo selbri lujvo lu se bridi li'u le bridi lu se bridi li'u .i&lt;br&gt; 117. zo tertau lujvo lu te tanru li'u le tanru lu te tanru li'u .i&lt;br&gt; 118. lu gismu smuni li'u tanru zo gismu zo smuni .i&lt;br&gt; 119. lu gismu lujvo li'u tanru zo gismu zo lujvo lu gismu co lujvo li'u .i&lt;br&gt; 120. lu le tanru li'u cu smuni lu le ve lujvo li'u mi .i&lt;br&gt; 121. zo tanru gismu zo'e le se tanru zoi ly. tau .ly .i&lt;br&gt; 122. le gismu sumti dei li pa .i&lt;br&gt; 123. nei sumti dei li re .i&lt;br&gt; 124. li ci sumti dei li ci .i&lt;br&gt; 125. zo sumti gismu .i zo sumti sumti di'u li pa .i&lt;br&gt; 126. mi klama le se klama le te klama .i le te klama sumti di'u li ci .i&lt;br&gt; 127. dei bridi le du'u bridi dei&lt;br&gt; 128. lu le ninmu klama le nanmu li'u bridi le du'u klama le nanmu&lt;br&gt; 129. do nanmu .i di'u bridi le du'u nanmu do .i&lt;br&gt; 130. di'e bridi le du'u cliva la'o gy. Elvis .gy .i la'o gy. Elvis .gy cliva .i&lt;br&gt; 131. zo .ui cmavo zoi selma'o. UI .selma'o lu mi gleki li'u la lojban. In Exercises 200 to 226, give an English translation of each place of each le sumti given.&lt;br&gt; 201. le smuni&lt;br&gt; 202. le se smuni&lt;br&gt; 203. le te smuni&lt;br&gt; 204. le gismu&lt;br&gt; 205. le se gismu&lt;br&gt; 206. le te gismu&lt;br&gt; 207. le ve gismu&lt;br&gt; 208. le lujvo&lt;br&gt; 209. le se lujvo&lt;br&gt; 210. le te lujvo&lt;br&gt; 211. le ve lujvo&lt;br&gt; 212. le tanru&lt;br&gt; 213. le se tanru&lt;br&gt; 214. le te tanru&lt;br&gt; 215. le ve tanru&lt;br&gt; 216. le xe tanru&lt;br&gt; 217. le sumti&lt;br&gt; 218. le se sumti&lt;br&gt; 219. le te sumti&lt;br&gt; 220. le bridi&lt;br&gt; 221. le se bridi&lt;br&gt; 222. le te bridi&lt;br&gt; 223. le cmavo&lt;br&gt; 224. le se cmavo&lt;br&gt; 225. le te cmavo&lt;br&gt; 226. le ve cmavo In Exercises 301 to 321, use the best word in the vocabulary of this lesson to fill in each blank of each sentence. You may use each word more than once.&lt;br&gt; 301. Teacher: _____ is the capital of Kentucky?&lt;br&gt; Pupil: Frankfort.&lt;br&gt; &lt;br&gt; 302. Mother: Did you hear what I just said?&lt;br&gt; Daughter: _____, except I missed the part about chocolate.&lt;br&gt; &lt;br&gt; 303. I said "_____", and she told me, "Sorry, I can listen right now."&lt;br&gt; 304. _____ name is Jessica, but my friends call me Jessie.&lt;br&gt; 305. Welcome to _____ .alaskas.!&lt;br&gt; 306. When I say "_____ Eliza", that means I'm talking to you.&lt;br&gt; 307. I wish I could stay longer. _____&lt;br&gt; 308. If I knew le _____, I wouldn't have asked you for a definition.&lt;br&gt; 309. Say "doi" and then le _____ of the person you wish to talk to.&lt;br&gt; 310. _____ Robert! It's so nice to see you after all these years.&lt;br&gt; 311. I wouldn't have said "_____" if I wasn't all ears.&lt;br&gt; 312. Did you rent the movie, "_____, Myself &amp; Irene"?&lt;br&gt; 313. Which le rafsi did you use to make that le _____?&lt;br&gt; 314. How can you call that le _____, when a phrase compound needs at least two words?&lt;br&gt; 315. If f(x) = y, then f is the function and yby. _____ fy.&lt;br&gt; 316. That's na ____ because any decent predicate needs at least pa le selbri.&lt;br&gt; 317. ro le attitudinals are a type of word called le _____.&lt;br&gt; 318. _____, but I would feel happy to listen another time.&lt;br&gt; 319. ____ la romeos. _____ la romeos. Wherefore art thou Romeo?&lt;br&gt; 320. _____ blue begins with the letter "b".&lt;br&gt; 321. Lojban has 1342 le _____. Answers to exercises. 1. Hello Brown. May I call you Bob?&lt;br&gt; 2. May I have a word with you, Mary? Let's leave.&lt;br&gt; The Lojban means more "We leave." than the English. Maybe I could change that with an attitudinal meaning "suggestion" or {ko} and {mi} instead of {mi'o}.&lt;br&gt; 3. The first person says, "The affix compound qualifies as an affix compound." The second person says, "I understand."&lt;br&gt; 4. The first person says, "I qualify as a structure word.". The second person says, "Huh?".&lt;br&gt; 5. Goodbye, Edna. Come see me in the future.&lt;br&gt; 6. Daniel says, "May I have a word with you?". Claire says, "Go ahead.".&lt;br&gt; 7. I said, "May I speak with you, Rachel?" The man associated with Rachel said, "Rachel said, 'I can't listen right now.'"&lt;br&gt; 8. I said the word {bridi}. I did not say the word {tanru}.&lt;br&gt; 9. I said, "What modifies what in the metaphor 'meaning sort of name'?" Kevin said, "In the metaphor 'meaning sort of name', 'meaning' modifies 'name'."&lt;br&gt; 10. I go by the name Mark.&lt;br&gt; 11. People associated with me call me "Brian". 100. According to me, "the relationship of a sentence" means "the relationship of a sentence".&lt;br&gt; 101. I see "Hello" meaning {coi}.&lt;br&gt; 102. I see {co'o} as meaning "Goodbye".&lt;br&gt; 103. I accept the interpretation "I am Brian." for the sentence "I am Brian.".&lt;br&gt; 104. George Bush accepts the interpretation of Mr. President as himself.&lt;br&gt; 105. I recognize the interpretation of the affixes of a root word as the affixes of a word.&lt;br&gt; 106. The word {gismu} (meaning root word) has the affix {gi'u}.&lt;br&gt; 107. The word {smuni} qualifies as a root word.&lt;br&gt; 108. The relationship of the word {gismu} has a role for the roles of the relationship of a root word.&lt;br&gt; 109. The word {selbri} denotes the relationship {se bridi}.&lt;br&gt; 110. The affix compound {selbri} means {se bridi}.&lt;br&gt; 111. The relationship of the compound word {selbri} has the same roles as the compound word {bridi}.&lt;br&gt; 112. The compound word {selbri} comes from the metaphor {se bridi}.&lt;br&gt; 113. The word {selma'o} comes from the metaphor {se cmavo}.&lt;br&gt; 114. The relationship of the word {selma'o} involves the grammatical class of a structure word.&lt;br&gt; 115. The compound word {selma'o} means {se cmavo}, and its relationship involves its meaning, and it comes from the metaphor {se cmavo}.&lt;br&gt; 116. The compound word {selbri} means {se bridi}, and its relationship involves the predicate, and it comes from the metaphor {se bridi}.&lt;br&gt; 117. The compound word {tertau} means {te tanru}, and its relationship involves a binary metaphor, and it comes from the metaphor {te tanru}.&lt;br&gt; 118. "Root word sort of meaning" qualifies a binary metaphor with "root word" modifying "meaning".&lt;br&gt; 119. In binary metaphor "Root word sort of compound word", "root word" modifies "compound word", and it means "compound word of root word sort".&lt;br&gt; 120. According to me, "metaphor made into compound word" means "binary metaphor".&lt;br&gt; 121. The root word {tanru} denotes the relationship involving the modifier of the metaphor, and has the affix {tau}.&lt;br&gt; 122. The root word is in the current utterance filling the first role.&lt;br&gt; 123. The current utterance is in the current predicate filling the second role.&lt;br&gt; 124. Three is in the current utterance filling the third role.&lt;br&gt; 125. The word {sumti} qualifies as a root word. The word {sumti} is in the previous utterance filling the first role.&lt;br&gt; 126. I go to the destination from the origin. The destination is in the previous utterance filling the third place.&lt;br&gt; 127. The current utterance qualifies as a predicate with the predicate relationship and the role the current utterance.&lt;br&gt; 128. The predicate "The woman goes to the man." involves the going relationship and the role of the man.&lt;br&gt; 129. You qualify as a man. The last utterance involves the man relationship and the argument you.&lt;br&gt; 130. The next utterance qualifies as a predicate involving the leaving relationship and Elvis. Elvis leaves.&lt;br&gt; 131. The structure word {.ui} comes from the grammatical class {UI}, means I feel happiness, and is used in the Lojban language. 201. meaning&lt;br&gt; 202. symbol with meaning&lt;br&gt; 203. interpreter of meaning&lt;br&gt; 204. root word&lt;br&gt; 205. relationship expressed by root word&lt;br&gt; 206. roles in relationship expressed by root word&lt;br&gt; 207. affixes of root word&lt;br&gt; 208. compound word&lt;br&gt; 209. meaning of compound word&lt;br&gt; 210. roles of relationship of compound word&lt;br&gt; 211. metaphor used to make compound word&lt;br&gt; 212. binary metaphor&lt;br&gt; 213. modifier of metaphor&lt;br&gt; 214. modified of metaphor&lt;br&gt; 215. meaning of metaphor&lt;br&gt; 216. usage of metaphor&lt;br&gt; 217. argument&lt;br&gt; 218. predicate of argument&lt;br&gt; 219. role filled by argument&lt;br&gt; 220. predicate&lt;br&gt; 221. relationship of predicate&lt;br&gt; 222. argument of predicate&lt;br&gt; 223. structure word&lt;br&gt; 224. grammatical class of structure word&lt;br&gt; 225. meaning of structure word&lt;br&gt; 226. language of structure word 301. ma&lt;br&gt; 302. je'e&lt;br&gt; 303. be'e&lt;br&gt; 304. mi&lt;br&gt; 305. la&lt;br&gt; 306. doi&lt;br&gt; 307. co'o&lt;br&gt; 308. smuni&lt;br&gt; 309. cmene&lt;br&gt; 310. coi&lt;br&gt; 311. re'i&lt;br&gt; 312. mi&lt;br&gt; 313. lujvo&lt;br&gt; 314. tanru&lt;br&gt; 315. sumti&lt;br&gt; 316. bridi&lt;br&gt; 317. cmavo&lt;br&gt; 318. re'inai&lt;br&gt; 319. doi, doi&lt;br&gt; 320. zo&lt;br&gt; 321. gismu 

Sound is defined as mechanical sinosodial vibratory longitudinal impulse waves which oscillate the pressure of a transmitting medium by means of adiabatic compression and decompression consequently resulting in the increase in the angular momentum and hence rotational kinetic energy of the particles present within the transmitting medium producing frequencies audible within hearing range, that is between the threshold of audibility and the threshold of pain on a Fletchford Munson equal loudness contour diagram. Intro. When two glasses collide, we hear a sound. When we pluck a guitar string, we hear a sound. Different sounds are generated from different sources. Generally speaking, the collision of two objects results in a sound. Sound does not exist in a vacuum; it travels through the materials of a medium. Sound is a longitudinal wave in which the mechanical vibration constituting the wave occurs along the direction of the wave's propagation. The velocity of sound waves depends on the temperature and the pressure of the medium. For example, sound travels at different speeds in air and water. We can therefore define sound as a mechanical disturbance produced by the collision of two or more physical quantities from a state of equilibrium that propagates through an elastic material medium. =Sound= The amplitude is the magnitude of sound pressure change within a sound wave. Sound amplitude can be measured in pascals (Pa), though its more common to refer to the "sound (pressure) level" as Sound intensity(dB,dBSPL,dB(SPL)), and the "perceived sound level" as Loudness(dBA, dB(A)). Sound intensity is flow of sound energy per unit time through a fixed area. It has units of watts per square meter. The reference Intensity is defined as the minimum Intensity that is audible to the human ear, it is equal to 10-12 W/m2, or one picowatt per square meter. When the intensity is quoted in decibels this reference value is used. Loudness is sound intensity altered according to the frequency response of the human ear and is measured in a unit called the A-weighted decibel (dB(A), also used to be called phon). The Decibel. The decibel is not, as is commonly believed, the unit of sound. Sound is measured in terms of pressure. However, the decibel is used to express the pressure as very large variations of pressure are commonly encountered. The decibel is a dimensionless quantity and is used to express the ratio of one power quantity to another. The definition of the decibel is formula_1, where x is a squared quantity, ie pressure squared, volts squared etc. The decibel is useful to define relative changes. For instance, the required sound decrease for new cars might be 3 dB, this means, compared to the old car the new car must be 3 dB quieter. The absolute level of the car, in this case, does not matter. Definition of terms. "Sample equation:" Change in sound intensity&lt;br&gt; Δβ = β2 - β1&lt;br&gt;  = 10 log("I"2/"I"0) - 10 log("I"1/"I"0)&lt;br&gt;  = 10 [log("I"2/"I"0) - log("I"1/"I"0)]&lt;br&gt;  = 10 log[("I"2/"I"0)/("I"1/"I"0)]&lt;br&gt;  = 10 log("I"2/"I"1)&lt;br&gt; where log is the base-10 logarithm. Doppler effect. &lt;br&gt;f' is the observed frequency, f is the actual frequency, v is the speed of sound (formula_2), T is temperature in degrees Celsius formula_3 is the speed of the observer, and formula_4 is the speed of the source. If the observer is approaching the source, use the top operator (the +) in the numerator, and if the source is approaching the observer, use the top operator (the -) in the denominator. If the observer is moving away from the source, use the bottom operator (the -) in the numerator, and if the source is moving away from the observer, use the bottom operator (the +) in the denominator. Example problems. A. An ambulance, which is emitting a 400 Hz siren, is moving at a speed of 30 m/s towards a stationary observer. The speed of sound in this case is 339 m/s. formula_5 B. An M551 Sheridan, moving at 10 m/s is following a Renault FT-17 which is moving in the same direction at 5 m/s and emitting a 30 Hz tone. The speed of sound in this case is 342 m/s. formula_6 

Our first review of the is in, by email to the author: Thanks karl! &lt;br&gt; it's very helpful!!!! 

Torque and Circular Motion. Circular motion is the motion of a particle at a set distance (called radius) from a point. For circular motion, there needs to be a force that makes the particle turn. This force is called the 'centripetal force.' Please note that the centripetal force is "not" a new type of force-it is just a force causing rotational motion. To make this clearer, let us study the following examples: Thus, we see that the centripetal force acting on a body is always provided by some other type of force -- centripetal force, thus, is simply a name to indicate the force that provides this circular motion. This centripetal force is "always" acting inward toward the center. You will know this if you swing an object in a circular motion. If you notice carefully, you will see that you have to continuously pull inward. We know that an opposite force should exist for this centripetal force(by Newton's 3rd Law of Motion). This is the centrifugal force, which exists only if we study the body from a non-inertial frame of reference(an accelerating frame of reference, such as in circular motion). This is a so-called 'pseudo-force', which is used to make the Newton's law applicable to the person who is inside a non-inertial frame. e.g. If a driver suddenly turns the car to the left, you go towards the right side of the car because of centrifugal force. The centrifugal force is equal and opposite to the centripetal force. It is caused due to inertia of a body. Average angular velocity is equal to one-half of the sum of initial and final angular velocities assuming constant acceleration, and is also equal to the angle gone through divided by the time taken. Angular acceleration is equal to change in angular velocity divided by time taken. Angular momentum. Angular momentum of an object revolving around an external axis formula_3 is equal to the cross-product of the position vector with respect to formula_3 and its linear momentum. Angular momentum of a rotating object is equal to the moment of inertia times angular velocity. Rotational Kinetic Energy is equal to one-half of the product of moment of inertia and the angular velocity squared. The equations for rotational motion are analogous to those for linear motion-just look at those listed above. When studying rotational dynamics, remember: 

Vocabulary 1. "See Vocabulary 1." Vocabulary 2. "See Vocabulary 2." 

cmavo.  le LE the described  non-veridical descriptor: the one(s) described as ...  cu CU selbri separator  elidable marker: separates selbri from preceding  sumti, allows preceding terminator elision  ti KOhA6 this here  pro-sumti: this here; immediate demonstrative it;  indicated thing/place near speaker  ta KOhA6 that there  pro-sumti: that there; nearby demonstrative it;  indicated thing/place near listener  tu KOhA6 that yonder  pro-sumti: that yonder; distant demonstrative it;  indicated thing far from speaker&amp;listener  goi GOI pro-sumti assign  sumti assignment; used to define/assign ko'a/fo'a  series pro-sumti; Latin 'sive'  ko'a KOhA4 it-1  pro-sumti: he/she/it/they #1 (specified by goi)  zo'o UI5 humorously  attitudinal modifier: humorously - dully - seriously  cf. xajmi, junri)  zo'onai UI*5 seriously  attitudinal modifier: humorously - dully - seriously gismu.  broda rod predicate var 1  1st assignable variable predicate (context determines  place structure)  (cf. cmavo list bu'a)  blanu bla blue  x1 is blue [color adjective]  (cf. skari, blabi, xekri, zirpu, kandi, carmi, cicna)  cusku cus sku express  x1 (agent) expresses/says x2 (sedu'u/text/lu'e  concept) for audience x3 via expressive medium x4  [also says]; (cf. bacru, tavla, casnu, spuda, cmavo  list cu'u, bangu, dapma, jufra, pinka)  ciska ci'a write ' scribe '  x1 inscribes/writes x2 on display/storage medium  x3 with writing implement x4; x1 is a scribe  [also x3 writing surface]; (cf. papri, penbi,  pinsi, tcidu, xatra, pixra, prina, finti for  'author' or specific authorial works, barna, pinka)  prenu pre person  x1 is a person/people (noun) [not necessarily human];  x1 displays personality/a persona  (cf. nanmu, ninmu, remna, zukte, sevzi)  remna rem re'a human  x1 is a human/human being/man (non-specific  gender-free sense); (adjective:) x1 is human  (cf. nanmu, ninmu, prenu)  nanmu nau man  x1 is a man/men; x1 is a male humanoid person [not  necessarily adult]  [word dispreferred in metaphor/example as sexist;  (use remna or prenu)]; (cf. ninmu, remna, prenu,  makcu, nanla, bersa)  ninmu nim ni'u woman ' women '  x1 is a woman/women; x1 is a female humanoid person  [not necessarily adult]  [word dispreferred in metaphor/example as sexist;  (use remna or prenu)]; (cf. nanmu, remna, prenu,  makcu, nixli)  gasnu gau do  x1 [person/agent] is an agentive cause of event x2;  x1 does/brings about x2  (cf. cmavo list gau, gunka, zukte, rinka, fasnu for  non-agentive events, jibri, kakne, pilno)  zukte zuk zu'e act  x1 is a volitional entity employing means/taking  action x2 for purpose/goal x3/to end x3  [also acting at, undertaking, doing; agentive cause  with volition/purpose; also x3 objective, end];  (cf. cmavo list zu'e, bapli, gunka, jalge, krinu,  mukti, rinka, snuti, gasnu, fasnu, minji, prenu,  ciksi, jibri, pilno, pluta, tadji, tutci)  rinka rik ri'a cause  x1 (event/state) effects/physically causes effect x2  (event/state) under conditions x3  [x1 is a material condition for x2; x1 gives rise to  x2]; (cf. gasnu, krinu, nibli, te zukte, se jalge,  bapli, jitro, cmavo list ri'a, mukti, ciksi, xruti) 

French accents on computers. While French keyboards are available, some French students may need to enter accented characters on an English keyboard. If you are on the Internet, many sites have a virtual keyboard that allows you to mouse-select the characters. Google Translate, for example, has a virtual keyboard icon for entering text in the form window. Windows Operating System. If you are using Microsoft Windows Operating System, then you can use the Character Map application, located under: Some word processing programs allow the user to enter accents using a key combination, while others may require an Alt code. The ALT code is entered by holding down the ALT key, and enter the number (all digits given) using the keypad (only the keypad). In most applications, you will need the "numlock" turned on to avoid undesirable effects. Mac Operating System. If you are using the Mac Operating System, there is a simple system that can be used with the Option (⌥) key. Open the System Preferences application (found in your Applications folder by default and in the Apple menu in the upper-left corner) and navigate to "Language and Text" preferences. Under the Input Sources tab, select U. S. Extended. Now you can use the following key combinations with the Option (⌥) key to form French accents. For instance: Press and hold ⌥ and then press e. Now you have a floating acute accent. Press e again to put the accent over that letter and form é. In the same way ⌥ - 6, then o will give you ô, etc. These shortcuts work throughout the operating system and do not depend on the application in which they are used. Linux Operating System. If you are using Ubuntu Linux with Gnome you select the Compose key from System: Preferences: Keyboard then Layouts: Layout Options: Compose key position. You can select one of Right Alt key, Left Win-key, Right Win-key, Menu key, Right Ctrl key or Caps Lock key (for a USA keyboard layout). The Keyboard preferences dialog has an area you can use to test the settings. See below for how to use the Compose key. Ubuntu with a different window manager, such as KDE should have a similar keyboard preferences utility. If you are using Unix or a derivative operating system (such as Linux) with XFree86, you can define a compose key by opening a terminal window and typing: 'To use the Windows menu key (between the right Windows key and right Ctrl key: 'To use the right Windows key: 'To use the right Alt key:' To use the Compose key, press and release the Compose key, then type two characters. Combinations useful for typing in French follow: 

=Waves= Wave is defined as the movement of any periodic motion like a spring, a pendulum, a water wave, an electric wave, a sound wave, a light wave, etc. Any periodic wave that has amplitude varied with time, phase sinusoidally can be expressed mathematically as Wave speed is equal to the frequency times the wavelength. It can be understood as how frequently a certain distance (the wavelength in this case) is traversed. Frequency is equal to speed divided by wavelength. Period is equal to the inverse of frequency. Variables&lt;br&gt; Definition of terms Image here The wave’s extremes, its peaks and valleys, are called antinodes. At the middle of the wave are points that do not move, called nodes. "Examples of waves:" Water waves, sound waves, light waves, seismic waves, shock waves, electromagnetic waves … Oscillation. A wave is said to oscillate, which means to move back and forth in a regular, repeating way. This fluctuation can be between extremes of position, force, or quantity. Different types of waves have different types of oscillations. Longitudinal waves: Oscillation is parallel to the direction of the wave. Examples: sound waves, waves in a spring. Transverse waves: Oscillation is perpendicular to direction of the wave. Example: light Interference. When waves overlap each other it is called interference. This is divided into constructive and destructive interference. Constructive interference: the waves line up perfectly and add to each others’ strength. Destructive interference: the two waves cancel each other out, resulting in no wave.This happens when angle between them is 180degrees. Resonance. In real life, waves usually give a mishmash of constructive and destructive interference and quickly die out. However, at certain wavelengths standing waves form, resulting in resonance. These are waves that bounce back into themselves in a strengthening way, reaching maximum amplitude. "Resonance is a special case of forced vibration when the frequency of the impressed periodic force is equal to the natural frequency of the body so that it vibrates with increased amplitude, spontaneously." 

Personal Pronoun Tables: nominative, genitive, dative &amp; accusative cases. Nominative case personal pronouns. The nominative case is used as the subject of a verb. Genitive case personal pronouns. The genitive case corresponds to the possessive case in English or to the English objective case preceded by 'of' and denoting possession. The use of genitive personal pronouns is very rare in German and many Germans are unable to use them correctly. Examples: Dative case personal pronouns. The personal pronouns in the dative case are used as indirect objects of verbs and after the prepositions aus, außer, bei, mit, nach, seit, von, zu. Accusative case personal pronouns. The personal pronouns in the accusative case are used as direct objects of transitve verbs and after the prepositions durch, für, gegen, ohne, um. 

=Standing waves= Wave speed is equal to the square root of tension divided by the linear density of the string. Linear density of the string is equal to the mass divided by the length of the string. The fundamental wavelength is equal to two times the length of the string. Variables&lt;br&gt; Definition of terms Fundamental frequency: the frequency when the wavelength is the longest allowed, this gives us the lowest sound that we can get from the system. In a string, the length of the string is half of the largest wavelength that can create a standing wave, called its fundamental wavelength. 

=Wave overtones= For resonance in a taut string, the first harmonic is determined for a wave form with one antinode and two nodes. That is, the two ends of the string are nodes because they do not vibrate while the middle of the string is an antinode because it experiences the greatest change in amplitude. This means that one half of a full wavelength is represented by the length of the resonating structure. The frequency of the first harmonic is equal to wave speed divided by twice the length of the string. (Recall that wave speed is equal to wavelength times frequency.) The wavelength of the first harmonic is equal to double the length of the string. The "nth" wavelength is equal to the fundamental wavelength divided by n. Harmonics for a taut string* Definition of terms The first overtone is the first "allowed" harmonic above the fundamental frequency ("F"1). In the case of a system with two different ends (as in the case of a tube open at one end), the closed end is a node and the open end is an antinode. The first resonant frequency has only a quarter of a wave in the tube. This means that the first harmonic is characterized by a wavelength four times the length of the tube. The wavelength of the first harmonic is equal to four times the length of the string. The "nth" wavelength is equal to the fundamental wavelength divided by n. Note that "n" must be odd in this case as only odd harmonics will resonate in this situation. Harmonics for a system with two different ends* &lt;br&gt; Vs: "velocity of sound"&lt;br&gt; 

Astronomy is the scientific study of celestial bodies in the visible universe, from the scale of a few meters to the macro scale, including: the underlying physics governing those bodies, what they are made of, their properties, distribution, relation, distance, movement, creation, age and demise. Our understanding of the universe has dramatically improved due to the progress of technology. Astronomy has been one of the most modernized areas of scientific study, but it is also one of the oldest sciences — practiced by all ancient civilizations to some degree. Sadly, people have increasingly started to lose connection with the observable universe, something that was previously even required for measuring time, and defining seasons. Astronomy is among our species' first technological steps, but today only passingly remarked about when it verifies something thought about in theoretical physics. Even in a highly industrialized global civilization, defined by consumerism only a few of us had the chance to go beyond simple images and concepts and have in considering what is around our blue dot and its implications for us. This Wikibook introduces the reader to that tapestry and the process that revealed it to humanity. It presents astronomy not only as a field of knowledge, but also as a human endeavor in science. Also on Wikimedia. Wikibooks is a Web site where this and other free textbooks are developed. In addition to this book, Wikibooks is the host of related textbook projects. 

General Astronomy &gt; The Solar System « Astronomy | The Sun » The Solar System may be broadly defined as that portion of the universe under the gravitational influence of the Sun. This includes the Sun itself as well as all planets, moons, asteroids, comets, dust, and ice orbiting the Sun. The Solar System is an example of a star system, which is similarly defined as that portion of the universe under the gravitational influence of one or more co-orbiting stars. The Solar System is a unitary star system, as it has only one star (Sol, our Sun). Components of the Solar System. The largest, most massive, and most prominent element of the Solar System is, of course, the Sun. The Sun makes up 99.8% of the mass of the Solar System. It is literally the point around which the entire Solar System turns. The Sun is virtually at the center of the Solar System; although gravity tugs by the planets may move the center of the System slightly away from the center of the Sun, it always resides deep within the Sun's core. The next largest objects in the Solar System are the planets. There are generally considered to be eight planets in the Solar System. They can be divided into two types: (1) the gas giant planets, which include Jupiter, Saturn, Uranus, and Neptune and (2) the terrestrial planets Mercury, Venus, Earth, and Mars. All eight planets orbit the Sun in elliptical, roughly circular orbits, in approximately the same plane. However, no planet orbits in exactly a circular orbit or exactly in the plane of the Sun's rotation. The orbit of Jupiter is the closest to the plane and circularity; the orbit of Pluto (a dwarf planet) deviates the most from both the plane and from circularity. After the eight major planets are the minor planets, asteroids and comets. Asteroids and comets are smaller objects than planets, but also orbit the Sun. Asteroids and comets are distinguished by their content: asteroids are primarily made up of rock, while comets are primarily made of ices and volatile compounds. Minor planets may be found anywhere in the Solar System, in orbits varying from circular to highly elliptical. Most, however, are found in three belts. The main asteroid belt is found between the planets Mars and Jupiter. As the name implies, it is made almost entirely of asteroids. The ../Kuiper Belt/ is found outside the orbit of Neptune, and encompasses the area from 30 to 100 astronomical units from the Sun. The Kuiper belt contains mainly comets, including very large comet-like objects called cubewanos or plutinos. Some astronomers also consider Pluto to be part of the Kuiper belt. The ../Oort Cloud/ is another belt of comets, and is believed to extend out to approximately one light-year from the Sun. Its existence is deduced from the frequent visitation of long-period comets, comets with extremely elliptical or even hyperbolic orbits. Arrangement of the Solar System. The Solar System may be divided by its components into three major regions: the inner system, the near outer system, and the far outer system. The near outer system might also be referred to as the middle system. The general term outer system refers to both the near and far outer systems. The inner system is composed of the Sun (the largest mass), the terrestrial planets (rocky and closer to the sun) and their moons(moons are drawn to planets because of their gravitational force), close-orbiting asteroids and comets, and the main asteroid belt. Objects in the inner system are almost exclusively composed of rock, with either no atmosphere or an atmosphere that composes little of the object's mass. The inner system's boundary is defined by the main asteroid belt, which separates it from the near outer system. The near outer system is composed of the gas giant planets and their moons, and asteroids and comets that orbit between the main asteroid belt and the Kuiper belt. Objects in the near outer system may have rock, liquid, gas, and ice as significant components. The near outer system's boundary is defined by the orbit of Neptune. The far outer system is composed of the ice planet Pluto, the Kuiper belt, the Oort Cloud, and comets that orbit between the belt and the cloud. Objects in the far outer system may have some rock components, but are mainly composed of ices. Boundary of the Solar System. The boundary of the Solar System is defined in two ways. The gravitational boundary may be described as the point at which objects no longer orbit the Sun. This boundary includes the Oort Cloud, but is poorly defined, as an object is not compelled to orbit the Sun at any point. Another definition is to declare the heliopause as the boundary of the Solar System. This boundary is more easily detectable and definable, but resides well within the Oort Cloud. « Astronomy | The Solar System | The Sun » | Stars, Clusters and Nebulae » 

Appendix 3 ~ Online Resources for German Language Students Online Wörterbücher - Dictionary. Slideshows with pictures and pronunciations. Language courses German at the time of insertion there is only one file about fruit - I will try to add new ones every week-end. Tandem. Tandem by E-Mail The Mixxer Tandem via Skype 

Introduction. This book aims to be a comprehensive source for any developer who is interested in programming for the Windows platform. It starts at the lowest level, with the Win32 API (C and VB Classic) and then goes over to MFC (C++). Beyond these basic sections, it will cover COM, and the creation of ActiveX modules from a variety of languages. Next, it delves into the Windows DDK, and talk about programming device drivers for Windows platform. Finally, it moves on to the highest-level programming tasks, including shell extensions, shell scripting, and finally ASP and WSH. Other topics that will be discussed here are: Writing screen-savers, creating HTML help modules, and compiling DLL files. This book will focus on topics that are specific to Windows, and avoids general programming topics. For related material the reader is encouraged to look into Wikibooks other works, they will cover general programming, ASM, C, C++, Visual Basic and Visual Basic.NET and other languages and concepts in greater detail. Appropriate links to these books are provided. The reader is assumed to have a previous knowledge of the programming languages involved. Specifically, prior knowledge in C, C++, and Visual Basic is required for certain sections of this book. Further Reading. Wikimedia Resources. Programming Languages: Information about Windows: Related topics: 

Vectors are quantities that are characterized by having both a numerical quantity (called the "magnitude" and denoted as |"v"|) and a direction. Velocity is an example of a vector; it describes the time rated change in position with a numerical quantity (meters per second) as well as indicating the direction of movement. The definition of a vector is any quantity that adds according to the parallelogram law (there are some physical quantities that have magnitude and direction that are not vectors). Scalars are quantities in physics that have no direction. Mass is a scalar; it can describe the quantity of matter with units (kilograms) but does not describe any direction. Frequently Asked Questions about Vectors. When are scalar and vector compositions essentially the same? Answer: when multiple vectors are in same direction then we can just add the magnitudes.so, the scalar and vector composition will be same as we do not add the directions. What is a "dot-product"? (work when force not parallel to displacement). Answer: Let's take gravity as our force. If you jump out of an airplane and fall you will pick up speed. (for simplicity's sake, let's ignore air drag). To work out the kinetic energy at any point you simply multiply the "value" of the force caused by gravity by the "distance" moved in the direction of the force. For example, a 180 N boy falling a distance of 10 m will have 1800 J of extra kinetic energy. We say that the man has had 1800 J of work done on him by the force of gravity. Notice that energy is "not" a vector. It has a value but no direction. Gravity and displacement are vectors. They have a value plus a direction. (In this case, their directions are down and down respectively) The reason we can get a scalar energy from vectors gravity and displacement is because, in this case, they happen to point in the same direction. Gravity acts downwards and displacement is also downwards. When two vectors point in the same direction, you can get the scalar product by just multiplying the "value" of the two vectors together and ignoring the direction. But what happens if they don't point in the same direction? Consider a man walking up a hill. Obviously it takes energy to do this because you are going against the force of gravity. The steeper the hill, the more energy it takes every step to climb it. This is something we all know unless we live on a salt lake. In a situation like this we can still work out the work done. In the diagram, the green lines represent the displacement. To find out how much work "against" gravity the man does, we work out the "projection" of the displacement along the line of action of the force of gravity. In this case it's just the y component of the man's displacement. This is where the cos θ comes in. θ is merely the angle between the velocity vector and the force vector. If the two forces do not point in the same direction, you can still get the scalar product by multiplying the projection of one force in the direction of the other force. Thus: There is another method of defining the dot product which relies on components. What is a "cross-product"? (Force on a charged particle in a magnetic field). Answer: Suppose there is a charged particle moving in a constant magnetic field. According to the laws of electromagnetism, the particle is acted upon by a force called the Lorentz force. If this particle is moving from left to right at 30 m/s and the field is 30 Tesla pointing straight down perpendicular to the particle, the particle will actually curve in a circle spiraling out of the plane of the two with an acceleration of its charge in coulombs times 900 newtons per coulomb! This is because the calculation of the Lorentz force involves a cross-product.when cross product can replace the sin0 can take place during multiplication. A cross product can be calculated simply using the angle between the two vectors and your right hand. If the forces point parallel or 180° from each other, it's simple: the cross-product does not exist. If they are exactly perpendicular, the cross-product has a magnitude of the product of the two magnitudes. For all others in between however, the following formula is used: But if the result is a vector, then what is the direction? That too is fairly simple, utilizing a method called the "right-hand rule". The right-hand rule works as follows: Place your right-hand flat along the first of the two vectors with the palm facing the second vector and your thumb sticking out perpendicular to your hand. Then proceed to curl your hand towards the second vector. The direction that your thumb points is the direction that cross-product vector points! Though this definition is easy to explain visually it is slightly more complicated to calculate than the dot product. How to draw vectors that are in or out of the plane of the page (or board). Answer: Vectors in the plane of the page are drawn as arrows on the page. A vector that goes into the plane of the screen is typically drawn as circles with an inscribed X. A vector that comes out of the plane of the screen is typically drawn as circles with dots at their centers. The X is meant to represent the fletching on the back of an arrow or dart while the dot is meant to represent the tip of the arrow. 

Work is equal to the integral of force times velocity times displacement times time. An integral, if the force isn't constant. 

A linear transformation is an important concept in mathematics because many real world phenomena can be approximated by linear models. Unlike a linear function, a linear transformation works on vectors as well as numbers. Motivations and definitions. Say we have the vector formula_1 in formula_2, and we rotate it through 90 degrees, to obtain the vector formula_3. Another example instead of rotating a vector, we stretch it, so a vector formula_4 becomes formula_5, for example. formula_6 becomes formula_7 Or, if we look at the "projection" of one vector onto the "x" axis - extracting its "x" component - , e.g. from formula_6 we get formula_9 These examples are all an example of a "mapping" between two vectors, and are all linear transformations. If the rule transforming the matrix is called formula_10, we often write formula_11 for the mapping of the vector formula_4 by the rule formula_10. formula_10 is often called the transformation. Note we do not always write brackets like when we write functions. However we "should" write brackets, especially when we want to express the mapping of the sum or the product or the combination of many vectors. Definitions. Linear Operators. Suppose one has a field K, and let x be an element of that field. Let O be a function taking values from K where O(x) is an element of a field J. Define O to be a linear form if and only if: Linear Forms. Suppose one has a vector space V, and let x be an element of that vector space. Let F be a function taking values from V where F(x) is an element of a field K. Define F to be a linear form if and only if: Linear Transformation. This time, instead of a field, let us consider functions from one vector space into another vector space. Let T be a function taking values from one vector space V where L(V) are elements of another vector space. Define L to be a linear transformation when it: Note that not all transformations are linear. Many simple transformations that are in the real world are also non-linear. Their study is more difficult, and will not be done here. For example, the transformation "S" (whose input and output are both vectors in R2) defined by formula_15 We can learn about nonlinear transformations by studying easier, linear ones. We often "describe" a transformation T in the following way This means that T, whatever transformation it may be, maps vectors in the vector space V to a vector in the vector space W. The actual transformation "could" be written, for instance, as Examples and proofs. Here are some examples of some linear transformations. At the same time, let's look at how we can prove that a transformation we may find is linear or not. Projection. Let us take the projection of vectors in R2 to vectors on the "x"-axis. Let's call this transformation T. We know that T maps vectors from R2 to R2, so we can say and we can then write the transformation itself as Clearly this is linear. ("Can you see why, without looking below?") Let's go through a proof that the conditions in the definitions are established. Scalar multiplication is preserved. We wish to show that for all vectors v and all scalars λ, T(λv)=λT(v). Let Then Now If we work out λT(v) and find it is the same vector, we have proved our result. This is the same vector as above, so under the transformation T, "scalar multiplication is preserved". Addition is preserved. We wish to show for all vectors x and y, T(x+y)=Tx+Ty. Let and Now Now if we can show Tx+Ty is this vector above, we have proved this result. Proceed, then, So we have that the transformation T "preserves addition". Zero vector is preserved. Clearly we have Conclusion. We have shown T preserves addition, scalar multiplication and the zero vector. So T must be linear. Disproof of linearity. When we want to "disprove" linearity - that is, to "prove" that a transformation is "not" linear, we need only find one counter-example. If we can find just one case in which the transformation does not preserve addition, scalar multiplication, or the zero vector, we can conclude that the transformation is not linear. For example, consider the transformation We suspect it is not linear. To prove it is not linear, take the vector then but so we can immediately say T is not linear because it doesn't preserve scalar multiplication. Problem set. Given the above, determine whether the following transformations are in fact linear or not. Write down each transformation in the form T:V -&gt; W, and identify V and W. (Answers follow to even-numbered questions): Images and kernels. We have some fundamental concepts underlying linear transformations, such as the "kernel" and the "image" of a linear transformation, which are analogous to the "zeros" and "range" of a function. Kernel. The "kernel" of a linear transformation T: V -&gt; W is the set of all vectors in V which are mapped to the zero vector in W, ie., Coincidentally because of the matr to the matrix equation Ax=0. The kernel of a transform T: V-&gt;W is always a subspace of V. The dimension of a transform or a matrix is called the "nullity".. Image. The "image" of a linear transformation T:V-&gt;W is the set of all vectors in W which were mapped from vectors in V. For example with the trivial mapping T:V-&gt;W such that Tx=0, the image would be 0. ("What would the kernel be?"). More formally, we say that the image of a transformation T:V-&gt;W is the set Isomorphism. A linear transformation T:V -&gt; W is an isomorphic transformation if it is: 

Chapter 1 Chapter 1. Introduction to Basic Ecology The objective of this wikibook, "A Study Guide to Basic Ecology", is to give the reader a better understanding of the way life functions on Earth, and how it is organized. This section will introduce basic concepts and definitions required to establish a course that is both biological and ecological in approach. You will learn that there are numerous terms in ecology which basically mean the same thing, and that there are also a plethora of views on precise definitions. As ecology emerged in the 19th and early 20th centuries, many concepts of ecological organization were proposed. These came from both zoologists and botanists—whom seldom read the other's literature—and from different schools of thought. Definition of Ecology. The term oekologie () was coined in 1866 by the German biologist, from the Greek "oikos" meaning "house" or "dwelling", and "logos" meaning "science" or "study". Thus, ecology is the "study of the household of nature". Haeckel intended it to encompass the study of an animal in relation to both the physical environment and other plants and animals with which it interacted. A contemporary definition of ecology is: The scientific study of the distribution and abundance of organisms and the interactions that determine distribution and abundance. This definition encompasses not only the plants and animals that Haeckel recognized but microscopic organisms such as , and , as well. The interactions that determine an organism's distribution and abundance are processes that include energy flow, growth, reproduction, predation, competition and many others. Basic Ecology Definitions. To fully understand the science of ecology, there are some common terms that must be defined. The term environment describes, in an unspecified way, the sum total of physical and biotic conditions that influence an organism (Kendeigh, 1961). The subset of the planet earth "environment" into which life penetrates is termed the biosphere. With respect to the planet earth, the biosphere penetrates only a limited distance into the rock beneath the land and the oceans, and a limited distance out away from the planet towards space. All human effort so far has failed to demonstrate that the biosphere extends beyond these limits or that other biospheres exist elsewhere in the universe. We cannot therefore conclude that they do not exist, only that we know nothing of that existence. Ecosystem is perhaps the most widely used term in ecology. It is defined as the system of organisms and physical factors under study or consideration. Although the boundaries of ecosystems are sometimes quite difficult to define in nature, ecosystems—however bounded—comprise the basic units of that nature (Tansley, 1935). Habitat is generally considered by biologists to be the physical conditions that surround a species, or species population, or assemblage of species, or community (Clements and Shelford, 1939). The basic physical units of the biosphere are the lithosphere (the land), hydrosphere (the water), and atmosphere (the air). Apparently there is no permanent biota of the atmosphere, although insects and birds among others utilize that environment extensively (Hesse, et al., 1951). These basic units are easily recognized in any landscape, as for example shown in the scene at right from Belarus in Europe. NOTE: Definitions of terms used in this textbook have been collected into various glossaries, the basic one being Basic Ecology Glossary. Scope of Ecology. Humanity has been studying nature for thousands of years and formally for several centuries under the science of biology, so why do we say that Ecology is a relatively new science? (Note, the term "ecology" is only about 150 years old). What is it that ecologists do or study that biologists might not? Kendeigh (1961, p. 3) wrote that "[e]cology is one of three main divisions of biology; the other two being morphology [organism structural aspects] and physiology [organism functional aspects]". The behavioral sciences encompass one facet of the interaction of an organism with its physical and biological surroundings, and therefore are part of ecology. Several techniques are used to study Ecology; one of the methods is known as strong inference. Strong inference is defined as the method of testing a hypothesis by deliberately attempting to demonstrate the falsity of the hypothesis (Kinraide and Denison, 2003). Performing experiments is the preferred method of using strong inference. The point of an experiment is to control variables and/or patterns and to predict the outcomes. Traditionally, Ecology is based more on weak inference (the other type of inference) rather than strong inference. In weak inference there are no experiments; rather, there are correlations between observations. Weak inference is defined as a "common sense" approach using simple system models to build understanding (Elner and Vadas). Both weak and strong inference have been used throughout many different science courses over the last several thousands of years (Kinraide and Denison). However, strong inference is said to be more advantageous because weak inference is more prone to error. Weak inference makes people try to cling to one point of view; whereas strong inference (doing experimentation) causes people to think from more than one point of view and keep coming up with different hypotheses to test as their old hypotheses get falsified. Ecology incorporates and overlaps with many other disciplines in both the biological and physical sciences. Certainly on one level, there is no information about the natural environment that does not have some applicability to ecology. On the other hand, such interrelationships between the sciences are much more the norm than most people realize. Ecology, according to Kendeigh (1961, p. 1), is a distinct science because of the way information about the natural world is organized; and because of its unique methodologies and point-of-view by which knowledge is extracted from such information. It may be helpful here to list the phenomenon that ecologists are particularly interested in if only to better clarify the boundaries of this text covering basic ecology (list modified from Kendeigh, 1961, p 1). Ecologists study: Ecology is both a biological and an environmental science, something that should certainly be evident from the definition provided above. Many environmental sciences are minimally concerned with biology (meteorology, for example) and others (environmental toxicology, for example) necessarily combine physical and biological sciences. Environmentalism on the other hand is a political position involved with various aspects of managing the environment. There are many subcategories of Ecology. Plant Ecology looks at the differences and similarities of various plants in differing climates and habitats. The origins of Animal Ecology can be traced to two Europeans, R. Hesse of and of . Physiological Ecology, or , studies the responses of the individual organism to the environment. The idea of was brought forward by two famous men, , who looked at the similarities and dissimilarites of populations and how they replaced each other over time, and . The last subcategory is , which brought about terms such as , introduced by A. Thienemann. Early History of Ecology. Because ecology is a facet of biology, it is not possible to clearly define a time or person that represents the beginning of ecological thought as distinct from biology. One of the first ecologists was probably Theophrastus, a contemporary and student of Aristotle. Theophrastus described various interrelationships between animals and between animals and their environment as early as the 4th century BC (Ramalay, 1940). Over the next several centuries biological knowledge gradually expanded as naturalists, such as Bufon and Linnaeus, contributed to a growing understanding of the nature of plants and animals. These early naturalists of course understood that organisms had a relationship with the environment, and especially that the distribution of forms was related to physical aspects of the geography. The modern ecological concept of integrated communities of organisms began in the 19th century with the studies of August Grisebach (1838), a German botanist. A contemporary and fellow countryman, Ernst Haeckel, coined the term "ecology" in 1866. Other biologists followed Grisbach's ecological approach to natural history studies: K. Möbius (1877) investigating Danish oyster banks, Stephen A. Forbes (1887), describing a lake community as a microcosm, J. E. B. Warming (1895), describing Danish plant communities, and C. C. Adams conducting ecological stuidies in northern Michigan (1905) and at Isle Royale (1909). References. Tansley, A. G. 1935. The use and abuse of vegetational concepts and terms. Ecology, 16(3): 284-307. 

Machines process things. We feed stuff into a machine and get different stuff out. A saw turns trees into planks. An internal combustion engine turns gasoline into rotational energy. A computer is no different. But instead of physical materials, computers process information for us. We feed information into the computer, tell the computer what do with it, and then get a result back out. The information we put into a computer is called input and the information we receive from the computer is called output. Input can come from just about anywhere. Keystrokes on a keyboard, data from an internet connection, or sound waves converted to electrical signals are examples of input. Output can also take many forms such as video played on a monitor, a string of text displayed in a terminal, or data we save onto a hard drive. The collection of input and generation of output is known under the general term, "input/output", or I/O for short, and is a core function of computers. Interestingly, the C programming language doesn't have I/O abilities built into it. It does, however, provide us with an external library containing I/O functions which we can compile and link into our programs. We have already used an output library function in the Hello, World! example at the beginning of this text: codice_1. You may recall this function resided in the codice_2 library file. As that file's name implies, codice_2 contains standardized I/O functions for adding input and output capability to our programs. This section of the text will explore some of these functions. Output using codice_1. Recall from the beginning of this text the demonstration program duplicated below: int main(void)  printf("Hello, World!");  return 0; If you compile and run this program, you will see the sentence below show up on your screen: This amazing accomplishment was achieved by using the "function" codice_1. A function is like a "black box" that does something for you without exposing the internals inside. We can write functions ourselves in C, but we will cover that later. You have seen that to use codice_1 one puts text, surrounded by quotes, in between the parentheses. We call the text surrounded by quotes a "literal string" (or just a "string"), and we call that string an "argument" to printf. As a note of explanation, it is sometimes convenient to include the open and closing parentheses after a function name to remind us that it is, indeed, a function. However usually when the name of the function we are talking about is understood, it is not necessary. As you can see in the example above, using codice_1 can be as simple as typing in some text, surrounded by double quotes (note that these are double quotes and not two single quotes). So, for example, you can print any string by placing it as an argument to the codice_1 function: And once it is contained in a proper codice_9 function, it will show: Printing numbers and escape sequences. Placeholder codes. The codice_1 function is a powerful function, and is probably the most-used function in C programs. For example, let us look at a problem. Say we want to calculate: 19 + 31. Let's use C to get the answer. We start writing  // can't be used without this header int main(void)  printf("19+31 is"); But here we are stuck! codice_1 only prints strings! Thankfully, printf has methods for printing numbers. What we do is put a "placeholder" format code in the string. We write:  printf("19+31 is %d", 19+31); The placeholder codice_12 literally "holds the place" for the actual number that is the result of adding 19 to 31. These placeholders are called format specifiers. Many other format specifiers work with codice_1. If we have a floating-point number, we can use codice_14 to print out a floating-point number, decimal point and all. Other format specifiers are: A complete listing of all the format specifiers for codice_1 is on Wikipedia. Tabs and newlines. What if, we want to achieve some output that will look like:  1905  312 + codice_1 will not put line breaks in at the end of each statement: we must do this ourselves. But how? What we can do is use the newline "escape character". An escape character is a special character that we can write but will do something special onscreen, such as make a beep, write a tab, and so on. To write a newline we write codice_24. All escape characters start with a backslash. So to achieve the output above, we write  printf(" 1905\n312 +\n-----\n"); or to be a bit clearer, we can break this long printf statement over several lines. So our program will be int main(void)  printf(" 1905\n");  printf("312 +\n");  printf("-----\n");  printf("%d", 1905+312);  return 0; There are other escape characters we can use. Another common one is to use codice_25 to write a tab. You can use codice_26 to ring the computer's bell, but you should not use this very much in your programs, as excessive use of sound is not very friendly to the user. Other output methods. codice_27. The codice_27 function is a very simple way to send a string to the screen when you have no placeholders or variables to be concerned about. It works very much like the codice_1 function we saw in the "Hello, World!" example:  puts("Print this string."); will print to the screen:  Print this string. followed by the newline character (as discussed above). (The codice_30 function appends a newline character to its output.) Input using codice_31. The codice_31 function is the input method equivalent to the codice_1 output function - simple yet powerful. In its simplest invocation, the scanf "format string" holds a single "placeholder" representing the type of value that will be entered by the user. These placeholders are mostly the same as the codice_1 function - codice_12 for integers, codice_14 for floats, and codice_37 for doubles. There is, however, one variation to codice_31 as compared to codice_1. The codice_31 function requires the memory address of the variable to which you want to save the input value. While "pointers" (variables storing memory addresses) can be used here, this is a concept that won't be approached until later in the text. Instead, the simple technique is to use the "address-of" operator, &amp;. For now it may be best to consider this "magic" before we discuss pointers. A typical application might be like this: int main(void)  int a;  printf("Please input an integer value: ");  scanf("%d", &amp;a);  printf("You entered: %d\n", a);  return 0; If you were to describe the effect of the codice_31 function call above, it might read as: "Read in an integer from the user and store it at the address of variable "a" ". If you are trying to input a "string" using "scanf", you should not include the &amp; operator. The code below will produce a runtime error and the program will likely crash.  scanf("%s", &amp;a); The correct usage would be:  scanf("%s", a); This is because, whenever you use a format specifier for a string (codice_20), the variable that you use to store the value will be an array and, the array names (in this case - a) themselves point out to their base address and hence, the address of operator is not required. Note that using codice_31 to collect keyboard input from the user can make your code vulnerable to Buffer overflow issues and lead to other undesirable behavior if you are not very careful. Consider using codice_44 instead of codice_31. Note on inputs: When data is typed at a keyboard, the information does not go straight to the program that is running. It is first stored in what is known as a buffer - a small amount of memory reserved for the input source. Sometimes there will be data left in the buffer when the program wants to read from the input source, and the codice_31 function will read this data instead of waiting for the user to type something. Some may suggest you use the function codice_47, which may work as desired on some computers, but isn't considered good practice, as you will see later. Doing this has the downfall that if you take your code to a different computer with a different compiler, your code may not work properly. 

The study of power series is aimed at investigating series which can approximate some function over a certain interval. Motivations. Elementary calculus (differentiation) is used to obtain information on a line which touches a curve at one point (i.e. a tangent). This is done by calculating the gradient, or slope of the curve, at a single point. However, this does not provide us with reliable information on the curve's actual "value" at given points in a wider interval. This is where the concept of power series becomes useful. An example. Consider the curve of formula_1 , about the point formula_2 . A naïve approximation would be the line formula_3 . However, for a more accurate approximation, observe that formula_4 looks like an inverted parabola around formula_2 - therefore, we might think about which parabola could approximate the shape of formula_4 near this point. This curve might well come to mind: In fact, this is the best estimate for formula_4 which uses polynomials of degree 2 (i.e. a highest term of formula_9) - but how do we know this is true? This is the study of power series: finding optimal approximations to functions using polynomials. Definition. A "power series" (in one variable) is a infinite series of the form or, equivalently, Radius of convergence. When using a power series as an alternative method of calculating a function's value, the equation can only be used to study formula_14 where the power series converges - this may happen for a finite range, or for all real numbers. The size of the interval (around its center) in which the power series converges to the function is known as the "radius of convergence". An example. this converges when formula_16 , the range formula_17 , so the radius of convergence - centered at 0 - is 1. It should also be observed that at the "extremities" of the radius, that is where formula_18 and formula_19 , the power series does not converge. Another example. Using the ratio test, this series converges when the ratio of successive terms is less than one: which is always true - therefore, this power series has an infinite radius of convergence. In effect, this means that the power series can "always" be used as a valid alternative to the original function, formula_24 . Abstraction. If we use the ratio test on an arbitrary power series, we find it converges when and diverges when The radius of convergence is therefore If this limit diverges to infinity, the series has an infinite radius of convergence. Differentiation and Integration. Within its radius of convergence, a power series can be differentiated and integrated term by term. Both the differential and the integral have the same radius of convergence as the original series. This allows us to sum exactly suitable power series. For example, This is a geometric series, which converges for formula_16 . Integrating both sides, we get which will also converge for formula_16 . When formula_19 this is the harmonic series, which "diverges"; when formula_18 this is an alternating series with diminishing terms, which "converges" to formula_36 - this is testing the extremities. It also lets us write series for integrals we cannot do exactly such as the error function: The left hand side can not be integrated exactly, but the right hand side can be. This gives us a series for the sum, which has an infinite radius of convergence, letting us approximate the integral as closely as we like. Note that this is not a power series, as the power of formula_39 is not the index. 

 Infinite Sums  Derivative Rules and the Substitution Rule  Integration by Parts  Trigonometric Substitutions  Trigonometric Integrals  Rational Functions by Partial Fraction Decomposition  Tangent Half Angle Substitution  Reduction Formula  Irrational Functions  Numerical Approximations 

Note that "quitter" must be followed by a direct object, usually a room or building. , meaning "to inhabit", "to dwell", or "to reside", is used to say in what city or area you live: "Habiter" is also used more specifically: "Habiter rue …" is used to state on what street a person lives: "Habiter" refers to occupying a location, and does not mean "to live" more generally. (The irregular "vivre" is used instead.) The verb "faire" is translated to "to do" or "to make". It is irregularly conjugated (it does not count as a regular "-re" verb). Examples. "Faire" conjugated, followed by an infinitive, means "to have something done for oneself":  The direct object pronouns "me", "te", "nous", and "vous" mean "me", "you", and "us": "Me", "te", "nous", and "vous" are also indirect object pronouns, and mean "to me", "to you", and "to us": These pronouns come before the verb they modify: "Me" becomes "m'" and "te" becomes "t'" before a vowel: Examples. The preposition means "in" or "into", in the sense of "inside", "from outside", or "to inside":  also means "in" in the sense of "within a period of time": As in English, describes abstract situations and state:  can also mean "out of" or "from": The preposition is used instead to indicate "in" in other senses. 

Romanian is a Romance language spoken mainly in Romania and Moldova, as well as in some parts of Hungary, Serbia, Bulgaria and Ukraine. It is part of the Romance group of the Indoeuropean family of languages. It is easiest to learn if someone already knows a related language such as Spanish, Catalan, French, Galician, Portuguese or Italian. The most closely related are the other Romance languages, Italian being the closest, but any knowledge of Spanish, French, Galician or Portuguese might be very useful, especially because of the lexical and vocabulary similarities and the grammatical structure. A bit of knowledge of the Latin grammar might be useful as well. Even more distant Indo-European languages have many similarities in both grammar and even common words, as many languages, like English have borrowed extensively words from Latin. It is useful to know the language if travelling in Romania, especially in rural areas. Many people who know English and Romanian will understand any of its relatives, especially Spanish, French or Italian. Note that in Romanian, there is a formal and informal form when addressing people. The informal is tu (you - singular) and voi (also you - plural). Use tu when addressing friends or people you know well. When addressing strangers, use Dumneavoastră or vă. Dumneavoastră is used for stressed pronouns, which are optional, "vă" is used for non-stressed pronouns, which are mandatory for direct and indirect objects. For instance, Vă multumesc means "Thank you" (accusative case, direct object), and Vă dau un telefon means "I telephone (to) you" (dative case, indirect object) The formal form of address requires the second person plural form of the verb at all times, even when addressing a single person. (This is similar to the French construction, and, to an extent, German.) Using the singular form can be considered rude or even insulting. 



There are seven different types of syllables in the English language. Schwa: This can end in a consonant or not, and is an unemphasized syllable whose vowel is somewhat swallowed and pronounced like "uh". Other syllable types can be reclassified as a schwa based on experience of how a word is regionally pronounced. In a dictionary, the schwa sound is written like an upside-down, lower case "e". Long and short vowels. Long vowels are the same sound as the name of the vowel, "a", "e", "i", "o", and "u". Short vowels are the hardest vowel sounds to pronounce in English. "ah" as in cat, "eh" as in pet, "ih" as in sit, "oo" as in not, "uh" as in nut. 

GCSE Science/Electricity So far we have looked at the effect of putting a current in a magnetic field and seeing that there is a force on the wire that carries the current. On this page we will look at a related effect known as "induction". The best way to learn about induction is to consider the results of a simple experiment. In fact if your school has a spot galvanometer, you should really take a look at these results "in the flesh". Otherwise you'll just have to take my word for it. Experimental set up to demonstrate induction. Look at the diagram below. A coil of wire (solenoid) is attached to the inputs of a spot galvanometer. The solenoid can be made by winding a piece of plastic coated wire a round a paper tube. A spot galvanometer (also known as a mirror galvanometer) looks like a fancy piece of kit, and it is pretty fancy, delicate and expensive, but it does a very simple job. It's really just a very sensitive ammeter. It measures current just like any other ammeter but is much more sensitive than the normal sort you usually see. It can detect "tiny" currents. The only other thing that is needed is an ordinary bar magnet. Experimental results. So, once the apparatus has been set up as in the diagram above, we can look and see what happens. Conclusions. Q1) What effect will reversing the magnet and the direction of movement have on any induced current ? Theoretic explanation of what's going on (Advanced). So now we know what happens we have to come up with an explanation. Look at the diagram below. In this diagram the green line represents a wire that is going into the screen or page. It is at right angles to the magnetic field between the two bar magnets. The is moved downwards so that it "cuts" the field lines. This induces an e.m.f. (voltage) across the length of wire. It is this e.m.f. that causes a current to flow if a complete circuit is made. So now we know why you get a current. It's caused by the induced voltage, but why do you get a voltage? Look at the diagram below. This is a close up of the wire seen end on. The blue circles represent where the magnetic field lines stick out of the screen or page. They are coming straight out of the screen. The wire is made of a metal and so has many free electrons. As the wire moves downwards each individual electron travels downwards too. They effectively form loads of little current flowing down. From the motor effect we know that if a current flows in a magnetic field there will be a force. in this case the force pushes the electrons to the right. It is the electrons all moving up to the right end of the wire that causes the voltage difference. So you see, Induction is really the motor effect. Q2) A student uses a stronger magnet. What effect do you think this will have on any induced current? Q3) Imagine more than one loop of wire cutting through the magnetic field. Each loop of wire has its own e.m.f. induced on it. If the wires are joined up, it's as if they are little cells in series. Fill in the blanks " As the number of windings on the solenoid "increases" the value of the induced voltage ______________. Q4) An engineer wants to uses the induction to create electrcity. He gets a couple of students to push an enormous magnet into and out of a huge coil (he uses a whip to make 'em go fast). In what way will the electricity that this set up produces be different to a the electricity produced by a battery? «The motor effect 

Transformers are devices that use induction to transform a high a.c. voltage to a low one or vice versa. They are incredibly important devices, and (even more importantly) they are examiners' favourites. On this page you will be looking at how a transformer works and how they are used. You will also be given some practice in using some simple equations that let you work out what the output voltage and current must be. How a transformer works. There is nothing new to learn in this section. You've already covered GCSE Science/Magnetic effects of a current and GCSE Science/Induction. This section used ideas from those two pages to understand a transformer. The simplest transformer consists of a soft iron core with two coils of wire independently wrapped around it. The first coil called the "primary coil" is connected to an a.c. power supply. The other coil, called the "secondary coil" is connected to the output. The current in the primary coil causes a magnetic field. It's just an electromagnet but because the current is a.c. the magnetic field is constantly changing. If we consider one end of the soft iron core, it's first a north pole, then the field gets weaker and weaker until it becomes zero, then it starts to build up and the end becomes a south pole. The process repeats and the end becomes a north pole again. Now let's think about what is going on in the secondary coil. The magnetic field of the first coil passes right through the secondary coil, what's more the field is constantly changing. This changing field induces a current in the secondary coil. So even though the coils are not physically connected, the current in one, causes the current through the other. If you imagine the lines of field they appear to pulse as the field grows stronger and weaker. When the field grows weak they sweep inwards cutting the secondary coil. Q1) What would the current in the secondary coil be like if you used d.c. in the primary coil? - d.c. in the primary coil will not create any current in the secondary coil. Transformers are used to change the voltage. If the number of turns in the secondary coil is greater than the number of turns in the primary coil, then the voltage in the secondary coil will be higher than that in the primary coil. This is called a "step-up transformer". Conversely, if the secondary coil has fewer turns than the primary coil, then the voltage will be lower, and it is called a "step-down transformer". A simple demonstration transformer. Your teacher may show you a demonstration of a transformer. (see here for details on how to set up a simple transformer). Look at the diagram below. The second coil is on a separate soft iron core. This core can be brought up to the primary coil. At large distances the bulb is off, but as the coils are brought closer together they suddenly attract one another (the primary coil is a powerful electromagnet). Once the two cores are stuck together the bulb in the secondary coil starts to glow. The brighness of the bulb depends on the number of turns on the secondary coil. The more windings there are, the brighter the bulb. Important equations. There are two things you need to remember about transformers. Step up and Step down transformers. Because the voltage in the secondary coil depends on the number of turns, we can use a transformer to "transform" one voltage to another. For example let's assume that the primary coil has 300 turns of wire and the secondary coil has 900 turns. This transformer will step up the voltage by a factor of 3. So if the primary coil has 2V applied to it, the secondary coil will have 6V across it. A step down transformer is the opposite of a step up transformer. Let's say the primary has 900 turns and the secondary has 450 turns. This time the voltage will be half on the secondary as it is on the primary. So if the primary has 2V applied, the secondary will have only 1V across it. Strategy for answering questions. Questions on transformers will often be of the form "A coil has 3000 turns of wire on a soft iron. Another coil of 12000 turns is brought into contact with the first coil. This second coil has a 12V bulb attached to it completing a circuit.What voltage should applied to the first coil?" to answer a question like this: Uses of transformers. The most important use of transformers is in the routing of mains electricity. The power disspated by a cable = I2R, where I is the current and R is the resistance. There is little that can be done to lower the resistance, so in order to lose as little energy as possible in heating up the cable, the current needs to be kept as low as possible. This is done using a step up transformer which increases the voltage to many thousands of volts. Because power is voltage times current, a high voltage means a low current for a given power value.(Notice that transormers do not obey Ohm's law). At the other end of the distribution network, the voltage is stepped down using another transformer so that the voltage going into homes and factories is not excessively high. 

Pronouncing Spanish based on the written word is much simpler than pronouncing English based on written English. This is because, with few exceptions, each letter in the Spanish alphabet represents a single sound, and even when there are many possible sounds, simple rules tell us which is the correct one. In contrast, many letters and letter combinations in English represent multiple sounds (such as the "ou" and "gh" in words like "cough", "rough", "through", "though", "plough", etc.). Letter-sound correspondences in Spanish. The table below presents letter-sound correspondences in the order of the traditional Spanish alphabet. (Refer to the article Spanish orthography in Wikipedia for details on the Spanish alphabet and alphabetization.) One letter, one sound. Pronouncing Spanish based on the written word is much simpler than pronouncing English based on written English. Each vowel represents only one sound. With some exceptions (such as "w" and "x"), each consonant also represents one sound. Many consonants sound very similar to their English counterparts. As the table indicates, the pronunciation of some consonants (such as "b") does vary with the position of the consonant in the word, whether it is between vowels or not, etc. This is entirely predictable, so it doesn't really represent a breaking of the "one letter, one sound" rule. The University of Iowa has a very visual and detailed explanation of the Spanish pronunciation. Here is another page with links to the audio files of the letters. Word stress. In Spanish there are two levels of stress when pronouncing a syllable: stressed and unstressed. To illustrate: in the English word ""thinking", "think" is pronounced with stronger stress than "ing". If both syllables are pronounced with the same stress, it sounds like "thin king"". With one category of exceptions ("-mente" adverbs), all Spanish words have one stressed syllable. If a word has an accent mark (´; explicit accent), the syllable with the accent mark is stressed and the other syllables are unstressed. If a word has no accent mark (implicit accent), the stressed syllable is predictable by rule (see below). If you don't put the stress on the correct syllable, the other person may have trouble understanding you. For example: "esta", which has an implicit accent in the letter "e", means "this (feminine)"; and "está", which has an explicit accent in the letter "a", means "is." "Inglés" means "English," but "ingles" means "groins." Adverbs ending in "-mente" are stressed in two places: on the syllable where the accent falls in the adjectival root and on the "men" of "-mente". For example: "estúpido" → "estúpidamente". The vowel of an unstressed syllable should be pronounced with its true value, as shown in the table above. Don't reduce unstressed vowels to neutral schwa sounds, as occurs in English. Rules for pronouncing the implicit accent. There are only the following rules for pronouncing the implicit accent. The stressed syllable is in bold letters: Any exception to these rules is marked by writing an acute accent ("máximo", "paréntesis", "útil", "acción"). In those exceptions, the stressed syllable is the one where the acute accent (called "tilde" in Spanish) appears. The diéresis ( ¨ ). In the clusters "gue" and "gui", the "u" is not pronounced; it serves simply to give the "g" a hard-"g" sound, like in the English word "gut" ("gue" → [ge]; "gui" → [gi]). However, if the "u" has a the diaeresis mark (¨), it is pronounced like an English "w" ("güe" → [gwe]; "güi" → [gwi]). This mark is rather rare. Examples: 

Even though Romanian spelling is very phonetic, a lot of foreigners find the language difficult. The accent and sounds are very similar to Italian (with slight Slavic influences), so remember to sound every letter clearly. Also, sounds very rarely differ between words (i.e. the letter s is always pronounced the same, everytime, unlike in English). In the pronunciation guide, X-SAMPA will be used to help those people who are familiar with it, or the IPA (International Phonetic Alphabet), in pronouncing Romanian. For those unfamiliar with it, the IPA is a phonetic alphabet designed to help in the description of foreign sounds, and is used extensively by linguists. A good place to start learning the IPA is here. X-SAMPA is an ASCII method of transcribing the IPA. A guide to it can be found here. 



Lernen 7-2 ~ Tour de France. (aus Wikipedia, der freien Enzyklopädie)&lt;br&gt; Die "Tour de France" ist eines der berühmtesten und wichtigsten sportlichen Großereignisse überhaupt. Seit 1903 wird die Tour alljährlich - mit Ausnahme der Zeit des Ersten und Zweiten Weltkriegs - drei Wochen lang im Juli ausgetragen und führt dabei in wechselnder Streckenführung quer durch Frankreich und das nahe Ausland. Eine Tour de France der Frauen ("grande boucle féminine") mit deutlich kürzeren Etappen wird seit 1984 gefahren. Sie steht medial völlig im Schatten ihres männlichen Pendants. Vokabeln 7A.  die Ausnahme exception  die Enzyklopädie encyclopedia  der Erste Weltkrieg WW I  das Großereignis major event  der Juli July  das Radrennen bicycle race  die Welt world  die Woche, die Wochen week, weeks  die Zeit time, period  der Zweite Weltkrieg WW II  (bei weitem) berühmteste among the most widely renowned, the most popular  alljährlich every year  bei among (one of)  berühmteste most celebrated, most renowned  frei, freien (Akkusativ) free  seit since  sportlich athletic  überhaupt altogether, generally  während during  drei Wochen lang three weeks lasting  weit broad, wide  wichtig important 



Puzzles | Decision puzzles | Yet another Weighing You have a balance and the following items of which you know the weights: a stapler weighing 80g, a pen weighing 40g, an eraser weighing 20g, a ruler weighing 10g, and a small coin weighing 5g. "How many different weights can you measure with these items?" solution 

Systematic Phonics If we can divide a word into syllables, it essentially divides it into easy, bite-size chunks that very easy to pronounce. In this way anyone can pronounce even very large and difficult words. There are a few simple steps to divide a word into syllables. Consonant patterns. Take a word and grab any two talking vowels (the two vowels should not have any other talking vowels in-between them). These are the bookends, the beginning and the end of the part of the word that we are going to look at. There are just five possible combinations of vowels and consonants that can occur: VV. This vowel pattern is always broken between the two vowels. That is, one syllable ends at the end of the first vowel, and another syllable begins with the beginning of the second vowel. VCV. This is the trickiest one of the five, the only one that will trip you up if you are not careful. First, tentatively divide the syllables between the first vowel and the consonant. Then pronounce the result. If it sounds right, this is the correct division. If it does not sound right, try dividing the syllables between the consonant and the second vowel. How do you know how to pronounce the syllables ? If the syllable ends with a consonant, it is a short vowel. If it ends in a vowel, it is a long vowel. Think about it this way: imagine that an ending consonant has hands, and can use them to shut things like a Jack-in-the-box. If the box is shut by that ending consonant, the Jack-in-the-box is stuffed down in the box and real short, just like the vowel will be short. If there is no consonant to shut the box, the Jack-in-the-box springs out and is real long, just like the vowel will be real long. VCCV. Always break syllables between the two consonants. Almost always, even if the two consonants are the same one repeated, like a double "s". The only exception would be if the two consonants are a single diagraph such as "th" and "sh". VCCCV. Look for a blend in the group of three consonants. Circle the blend. Then divide the syllables between the blend and the remaining consonant. VCCCCV. The four consonants in the middle can be made up of either two double blends or one triple blend with another consonant. Identify and circle the blend(s) involved, and divide between them. 

Systematic Phonics Blends or clusters are groups of consonants whose sounds blend together. They can be composed of two or three consonants and can begin or end a syllable. Examples include br, cr, dr, fr, gr, pr, tr, bl, cl, fl, gl, pl, sl, sm, sn, sp, st, str, sw, sc, sk, scr, squ, spl, spr, thr, tw, rl, ng, nt, nk, lk, mp, lt, nd, and rt. 

Systematic Phonics Systematic phonics is a toolbox designed to enable and empower a person to read. It is useful for children first learning to read, non-native speakers who are learning to read, and even the many native adult English speakers who have not yet learned to read. With systematic phonics, a reader can break down any word, even the most difficult, into maneagable pieces called syllables. When a person can correctly divide a word, he can then pronounce each part and then the whole word. Systematic phonics was developed decades ago by language experts as a way to decode the entire English language. Extensive studies support that the method is effective. However, systematic phonics has not been adopted by most educators, who tend to use "whole language" or "holistic learning" to teach kids to read. Proponents of systematic phonics claim that the failure of many students to read to level can be traced to the use of other, less systematic, teaching methods. Systematic phonics an be an important part of learning English as a Second Language. 



A great way to learn to do design is to stretch your graphic design muscles. That is, push up your shirtsleeves and get down to doing some projects. Don't worry if the stuff turns out bad or people hate it, it's all a process of learning. Business card design. Every professional needs a business card. The US standard is 2" x 3.5" (89 mm x 51 mm). Use a .125" (3.175 mm) margin on all sides. Using a desktop publishing program, such as Scribus, make a business card for yourself, with typography and graphic elements. Plan layout with an eye for hierarchy. Symbol design. Using Adobe Illustrator or a similar vector-graphics tool, design a stylized icon or symbol representing an animal, taking care to communicate the socially perceived emotional connotations of this animal. The icon should be scalable to the point that it can be used on office stationary as well as blown up to billboard size. The symbol should employ limited colors so as to be easily and inexpensively printable. Compact disc design. Find an existing or imagined compact disc recording and design a complete package for it, including all copy appropriate for a disc to be sold in a store (song list, copyright information, etc.) You may use photography that you take yourself or any drawings or other computer art, etc. that you create yourself or commission for the project. Design the inside and outside of the CD case as well as the disc itself, keeping yourself open to the possibilities of non-traditional CD cases (transparent, etc.) Typeface design. Design a complete alphabet with special characters of your own, original font. You may choose a serif or sans-serif font, or a decorative font. Give samples of bold, extra bold, and italic versions of some letters. Webpage design. Design, using appropriate tools, an entire webpage devoted to a social or political cause of your choosing. Use colors, images, and a style that conforms to the page's message. Importantly don't forget - that the page you are designing for the Visitors only. So first classify your visitors, and use the color schemes which they impress most. And try to design a user friendly navigation system with a low flashy images. Toy design. Design a novel toy, optimized for maximum "fun-ness" for an audience of your choosing, and its complete packaging. Include suggestions for a marketing campaign promoting the toy. Personal projects. Design a project that personally touches you. Whether it be a collage or a poster for something you believe in, it is important to perform your skills for your own use as well as for your clients. Combine tools (Adobe with Macromedia with ink, pencil, highlighters and your trusty scanner). Have fun. 

Biochemistry « Metabolism and energy It is difficult to discuss the subject of biochemistry without a firm foundation in general chemistry and organic chemistry, but if one doesn't remember the concept of the Acid Dissociation Constant (p"K"a) from Organic Chemistry, one can read up on the topic below. "Buffers" are essential to biochemical reactions, as they provide a (more or less) stable pH value for reactions to take place under constantly changing circumstances. The pH value in living cells tends to fall between 7.2 and 7.4, and this pH level is generally maintained by weak acids. (The pH values in lysosomes and peroxisomes differ from this value, as do the pH measurements of the stomach and other organs found in various types of plants and animals.) An "acid" is here defined simply as any molecule that can release a proton (H+) into a solution. Stronger acids are more likely to release a proton, due to their atomic and molecular properties. The tendency of an acid to release a proton is called the "dissociation constant" ("K"a) of that substance, with for HA &lt;-&gt; H+ + A-. A larger "K"a value means a greater tendency to dissociate a proton, and thus it means the substance is a stronger acid. The pH at which 50% of the protons are dissociated can therefore be calculated as: This equation is known as the Henderson-Hasselbalch equation. The Henderson-Hasselbalch equation is derived from the adjacent Ka expression. By taking the logarithm of base ten to both sides, the next part of the equation is obtained. Using the logarithmic property of multiplication, the [H+] breaks from the expression. Since log Ka is equal to -pKa and log [H+] is equal to -pH, they are then substituted. To obtain what is known as the Henderson-Hasselbalch equation, -pKa and -pH are subtracted from their respective sides to yield a positive equation. The Henderson-Hasselbalch equation interestingly enough predicts the behavior of buffer solutions. A solution of 1 M ethanoic (acetic) acid [HA] and 1 M sodium ethanoate [A-] will have a pH equal to the p"K"a of ethanoic acid: 4.76. If we added acid to pure water up to a concentration of 0.1 M, the pH would become 1. If we add the same concentration of acid to the buffer solution, it will react with the ethanoate to form ethanoic acid. The ethanoate concentration would drop to 0.9 M and the ethanoic acid concentration should rise to 1.1 M. The pH becomes 4.67 - very different from the pH=1 "without" a buffer. Similarly, adding 0.1 M alkali changes the pH of the buffer to 4.85, instead of pH=13 without buffer. Due to the amphipathic nature of amino acids - which are the monomer building blocks of all proteins, physiological conditions are always considered to be buffered, which plays a major role in the conformations and reactivities of substrates in the cell's liquid interior, its cytosol. A very small (which would include a large negative value) p"K"a indicates a very strong acid. A p"K"a value between 4 and 5 is the most common range for organic acid compounds. « Metabolism and energy 



=Drawing= It's a good idea for a student of graphic design to learn how to draw. Sure, you can bypass this and do everything by computer, but drawing by hand trains your mind and your eye to seek out details and improve your grasp of visual forms. It develops different parts of your mind that relate to composition, color, anatomy, depth and mood, and helps you directly perceive things and interpret them. Much of what we relate to in our lives is linear. For effective visual communication you need to develop a nonlinear, direct connection to the actual tensions, feelings, and dimesions you are interpreting and expressing in your work. Drawing can be an invaluable way to develop a direct connection to the world around you and develop your own outlook. The ultimate aim of the graphic designer is to understand, use and coordinate the tools of the craft (typeface, art, photography, printing, etc.) and create a message that rises above them. As in music, you do not pay attention to the instruments alone but to the composition as a whole. Draw some architecture, some life studies, and some abstract compositions. Start with basic media like pencil and charcoal, and then try some things in color. Look at master drawings for ideas and inspiration. In interpreting three dimensions on a flat surface with a defined page shape, you encounter a number of issues. Evaluate how light falls across an object and how you can make shapes out of shadows. Shading, the gradation of darkness, is one of the most important factors when you draw an object. Shading and defining basic forms can provide a truer understanding of a dimensional object on a flat surface, and can make the difference between a really shoddy piece and a proper drawing. Get beyond outlining in your definition of form. Appreciate your overall use of the page. Be familiar with positive and negative space (negative space is that space that falls between the objects) and understand they are all part of the overall composition. To be successful at graphic design you need to develop a taste for solving visual problems and a set of strategies for resolving them. 

=Typography= Typeforms are an integral part of modern communication. Since the invention of the printing press, people have used printed —and more recently— electronic type to communicate. Could you imagine a newspaper or magazine where all the articles were handwritten and copied? Imagine how difficult it would be to compose an e-mail or use desktop publishing software to create party invitations without the existence of typefaces. Typefaces are used in print and electronic media not just to communicate a legible message, but as an expressive tool of the author or designer. Each typeface can be used to convey a different style or atmosphere. There are typefaces evocative of Art Deco, The old American West, Weddings and other traditionally formal ceremonies, childhood, The Middle Ages, handwritten text, and an endless variety of other styles. It is almost imperative as a graphic designer, and "absolutely" imperative as a typographer, to develop an appreciation and understanding of both modern and historic typography. In everyday life, one should pay attention to the stylistic and practical uses of typefaces in various kinds of media. Note what feelings the designer was hoping to convey, or what style they were attempting to mimic. Classification. A typeface is a style of lettering, such as Helvetica or Times. A font is the set of a typeface, used to produce the letters. On a computer, it is a file used by the system. People often confuse "font" with "typeface". For example, Helvetica, point 12 is a different font from Helvetica, point 14, even though both are of the same typeface. A set of similar typefaces is called a "family." Within a family, typefaces are categorized as parent (e.g. Times, Helvetica) or relative (e.g. bold, italic). Typefaces are categorized also according to style (e.g. italic, book), weight (e.g. bold or light), and width (e.g. expanded). Type is measured in points, from top to bottom of the letters' invisible bounding boxes. Since each letter fills up a percentage of that bounding box, when the size of the box is increased, the letter size increases. The divisions you may already be familiar with are the serif and sans-serif fonts. Serif fonts, like Times, have the little feet and variable line width characters which make them easy to read. Sans-serif fonts, like Helvetica and Verdana, are drawn with more even-width lines and don't have the little feet, which gives them a clean, modern feel. Usage in Design. Text can have a number of different purposes in a design. It can be used for pure graphic appeal —or aesthetically— a case in which legibility may be less important than aesthetics. An example of graphical use of text is possibly that of major titles. Text may also, more commonly, be used for communication on a linguistic basis rather than visual, in which case, legibility is always the priority of the designer. A few examples of usage of text for linguistic communication are: the body text of an article or book, the text of a restaurant menu, or in product descriptions in a catalog. A font may be used either successfully or poorly, depending on its degree of relevance in the project and the skill of the designer. The designer must pay careful attention to letter, word, and line spacing as well as the size of the typeface and its stylistic contribution to the overall aesthetic of the project. He or she must optimize the properties of the text for its purpose in the overall design (aesthetic or communicative) and maintain legibility where it is necessary, and the designer is also expected to add visual variety with formatting and layout, as well as possible font changes where applicable. Choosing a Font. A designer has to choose a font that is not only appropriate to the mood of the design, but that is appropriate for the text's purpose in the design. For example, there are many kinds of decorative typefaces that one would not want to set an entire article in. This is because a purely decorative typeface tends to be distracting to the content of the message and tire the reader's eyes when used in large portions. Our eyes are most comfortable reading less idiosyncratic typefaces. Decorative typefaces are better suited for display type (greater than 14 points), while simpler type is better used for text (less than 14 points). To maintain readability in large blocks of text —such as in an article— stay consistent, and use only one family. Readability in Different Media. Standard graphic design wisdom holds that of the categories of serif and sans-serif, serif fonts are easier to read. This is because when reading, the eye quickly scans the tops of the letterforms, and a serif font has more immediately recognizable features thanks to the tiny 'eye-holds' provided by the serifs. There is a notion slowly gaining acceptance, that, for the purposes of purely electronic design, the reverse is sometimes true. Sans-serif fonts are more readable in this case because a screen has a lower resolution than a printed page, so the serifs only serve to smudge the letter forms. Some typefaces such as Times, originally designed for the London Times newspaper, or Futura, designed as a letterpress (raised plate) type for printing on paper, were intended for the printed page. Others such as Georgia and Verdana were designed for the lower resolution of text on a screen. The shapes of these typefaces are developed to optimize visibility in smaller sizes on a computer monitor. In larger sizes, these differences don't matter as much. "For maximum readability on the page: Serif; On the screen: Sans." Typeface Sizing. Keep the typeface at a reasonable size for reading. The numbered size of a typeface may reflect the overall height of the lines that stick out of the type, but not the readability size that relates to the inner dimensions of the letters. The type size should be chosen on a visual basis, and not purely on that of font size numbers. Usually, type that is proportioned so that the lower case size is larger in proportion to the overall height of the font, can produce a greater amount of legible words into an equivalent space. Typeface Spacing. Word spacing should be so that the reader is aware of the beginnings and endings of words with little or no difficulty. Experiment with spacing settings to find the best one. Letter spacing and line spacing may be used to expand on the expressivity of the font. Leave enough space between the lines so that the text is legible. Experimentation is important here, too. The reader's eyes should be pulled to the next word as they read, not to the lines above and below. The letter spacing should be that so the reader can easily differentiate between different characters. Yet again: experiment, experiment, experiment. When to Break Rules. With text used purely for graphical or display purposes, spacing, fonts, or colors that would be considered otherwise unreasonable can be utilized to create a visually appealing effect. The legibility rules that are extremely important in body text aren't as critical to effect the larger typeface size. Make sure that before deciding to use more untraditional methods of formatting text, you consider the desires of your client, and/or the stylistic effect the formatting will have on the overall design. Is it supposed to be a classic, elegant wedding invitation, or is this a layout for a skateboarding magazine article? Consider the circumstances and break or follow the rules of tradition appropriately. Remember that a common amateur mistake is to break too many —or the wrong— formatting rules in rebellion of the idea that aesthetics have constraints. Then, they end up with a horrible mess of a design that looks good under no circumstances. "The purpose of a graphic designer is to merge optimal form and function". Remember: anything can be art, but not anything can look good! Layout. A good layout is one that shows good use of the elements and principles of design. Most importantly, a designer should use the principles of design to draw the reader's eye both to and through the design easily. The elements of design are: color, value, texture, shape, form, space, and line. The principles of design are: contrast, emphasis, balance, unity, pattern, movement, and rhythm. The specifics of both the principles and elements vary from source to source, but the idea remains the same. If you haven't studied the principles and elements of design before, read more in the chapter Principles of Design, or scroll down to view the external links. =External Links= 

Like drawing, you can skip photography and still be a salary-earning graphic design professional. However, some familiarity with the craft will make you an even better designer. Kids today have many more options open to them than we did back when I graduated three years ago. That is, you can take digital pictures and save the money and time needed to develop the film. 

=Form and function= Just like in architecture, one of the fundamental debates in graphic design is between form and function. Form is defined as the shape or visual quality of something. It refers to aesthetics, how a piece looks. Much of graphic design centers on how to make a work appealing (or unappealing, or any other quality depending on the project goals). Function relates to getting the job done. Function is pragmatic and business-oriented. Printing newspaper on glossy magazine paper is too expensive, and printing it on tissue paper would fall apart. An artistic photograph of a mountain may not convey a message appropriate to an Arctic cruise. 

=Statistics= Looking at designers as a whole, only about 60 percent of entering designers stay in the field after two years. By five years, about 30 percent are left. Princeton Review 

This book was started by , a graduate of the acclaimed University of Cincinnati College of Design, Architecture, Art and Planning (DAAP). Note from the author. There are not many textbooks published on the discipline of graphic design. Perhaps this is due in part to the abstract, subjective nature of many of its aspects, and the fact that there are few universally-agreed-upon rules like there are for mathematics or the hard sciences. This book is an attempt to formalize some of the concepts that were taught to me while in design school. The material was difficult and sometimes even overwhelming for me at the time. I did not graduate at the top of my class. However DAAP is a widely respected school and I did receive the vigorous training that characterizes its graphic design department. In this book I will try to break down in a simple and understandable way the basics of what I learned. 

Congruency. "Congruent" shapes are the same size with corresponding lengths and angles equal. In other words, they are exactly the same size and shape. They will fit on top of each other perfectly. Therefore if you know the size and shape of one you know the size and shape of the others. For example: Each of the above shapes is congruent to each other. The only difference is in their orientation, or the way they are rotated. If you traced them onto paper and cut them out, you could see that they fit over each other exactly. Having done this, right away we can see that, though the angles correspond in size and position, the sides do not. Therefore it is proved the triangles are not congruent. Similarity. "Similar" shapes are like congruent shapes in that they must be the same shape, but they don't have to be the same size. Their corresponding angles are congruent and their corresponding sides are in proportion. 

GCSE Science/Electricity On the last page we looked at induction. That is how the motor effect acting on electrons in wires causes them to move and create a current. On this page we will look at some practical uses of induction and how electricity is made in power stations. Look at the diagram of an electric motor. As the coil rotates it cuts the field lines in a downwards direction then an upwards direction. This means that the electricity produced alternates. The current flows one way then the other. The frequency depends on the speed of rotation. Q1) A motor is rotated manually to produce a current. An oscilloscope is used to investigate this current and produces a trace like the one shown. The time is on the x axis and voltage is on the y axis. The motor is speeded up. Which one of the following shows the new trace? Alternating Current. In the section above you learned that a dynamo produces a.c. If you've ever looked closely at a school power pack you will probably have noticed that it has a.c. and d.c. outputs. D.c. stands for direct current. The current flows in one direction only and (apart from when it is switched on or off) doesn't change in value. A.c. constantly changes. An a.c. curve is typically that of a sine wave. The current switches back and forth and the value of the current is constantly changing.(This fact is important when we come to study transformers.) Some devices work equally well with either a.c. or d.c. Others must have the correct type or they won't work. For example think about a child's train on a track. If a.c. electricity is used the train will not go anywhere! This is because the current is 'telling' the motor to go forwards, go backwards, no, go forwards and so on. Anything with a small electric motor in it has to use d.c. (Special electric motors which use a.c. do exist. However, the ones you'll come across in school or in any toy application will almost certainly be d.c. motors.) On the other hand a bulb has no problem working on a.c. This is because a bulb converts the energy of the electric current to heat and light. It makes no difference which way the current flows, the energy is still there. So bulbs, electric fires, cookers and so on do not care if they use a.c. or d.c. Transformers (which we will study in the next section) will only work with a.c. This is because a transformer needs to have a constantly varying current, and d.c. only varies when it is switched on or off. What useful electrical component was the engineer thinking of?---- Answers | Transformers» 

Verbs are often called action words that show what the subject (a noun or pronoun) is doing. A verb is a word that signifies to be, to act, or to be acted on: as, "I am, I rule, I am ruled, I love, you love, he loves". Verbs are so called, from the Latin "verbum", a word; because the verb is that word which most essentially contains what is said in any clause or sentence. Although described as "action words", they can describe abstract concepts. They are a requirement of any sentence. Verbs have modifications of four kinds: moods, tenses, persons and numbers. Morphological forms. An English verb has four morphological forms (forms of word formation) ever needful to be ascertained in the first place: the present, the past, the present participle, and the past participle. The third person singular is the fifth morphological form. The present is that form of the verb, which is the root of all the rest; the verb itself; or that simple term which we should look for in a dictionary: as, "be, act, rule, love, defend, terminate". The past is that simple form of the verb, which denotes time past; and which is always connected with some noun or pronoun, denoting the subject of the assertion: as, "I was, I acted, I ruled, I loved, I defended". The present participle is that form of the verb, which ends commonly in "ing", and implies a continuance of the being, action, or passion: as, "being, acting, ruling, loving, defending, terminating". The past participle is that form of the verb, which ends commonly in "d" or "ed", and implies what has taken place: as, "been, acted, ruled, loved". Regularity. English, like many Germanic languages, contains both strong (or irregular, which is not "quite" the same as strong) and weak (regular) verbs. Irregular verbs are one of the most difficult aspects of learning English. Each irregular verb must be memorized, because they are not often easy to identify otherwise. Verbs are divided, with respect to their regularity, into four classes: regular and irregular, redundant and defective. A regular verb is a verb that forms the past and the past participle by assuming "d" or "ed": as, "love, loved, loving, loved". An irregular verb is a verb that does not form the past and the past participle by assuming "d" or "ed": as, "see, saw, seeing, seen". A redundant verb is a verb that forms the past or the past participle in two or more ways, and so as to be both regular and irregular: as, "thrive, thrived or throve, thriving, thrived or thriven". A defective verb is a verb that forms no participles, and is used in but few of the moods and tenses: as, "beware, ought, quoth". Persons and numbers. The person and number of a verb are those modifications in which it agrees with its subject. There are three persons and two numbers: thus, Where the verb is varied, the third person singular in the present tense, is regularly formed by adding "s" or "es": as, "I see, he sees; I give, he gives; I go, he goes; I fly, he flies; I vex, he vexes; I lose, he loses". Where the verb is not varied to denote its person and number, these properties are inferred from its subject: as, "if I love, if he love; if we love, if you love, if they love". Tenses. Tenses are those modifications of the verb, which distinguish time. There are three tenses - Each of the above category lists subcategories. One could even say there are twelve tenses because each of those comes in simple and in progressive forms, which have different meaning. The past tense is sometimes called imperfect, but the names perfect and imperfect do not fit their meaning. These names were derived from Latin where they were correct. The Present Simple Present Tense is that which expresses what now exists, is normal or correlated to senses. It is used with adverbs like "always", "generally". Present Continuous Tense is that which expresses what is temporary: Present perfect tense is that which expresses what has taken place, within some period of time not yet fully past, or is still valid. It is used with adverbs like "ever", "never", "today", "this week". Present perfect continuous tense is that which which started in the past and has not yet finished. The Past Simple Past tense is that which expresses what took place in time fully past. It is used with adverbs like "yesterday", " last week". Past continuous tense is that which expresses what was taking place when (suddenly) something else occurred. Past perfect tense is that which expresses what had taken place, at some past time mentioned, before something other happened. Past perfect continuous tense is that which expresses what had started before and was still going on, when something else occurred. The Future Simple Future Tense is that which expresses what will take place hereafter. Future continuous tense is that which expresses what will be currently taking place at a certain time in future. Future Perfect Tense is that which expresses what will have taken place at some future time mentioned. Future Perfect continuous Tense is that which expresses what will have started at some time and will still be ongoing, at some future time mentioned. Signification. An active verb is a verb in an active sentence, in which the subject performs the verb: as, An active verb can be transitive or intransitive, but not passive or neuter. Verbs are divided again, with respect to their signification, into four classes: transitive, intransitive, passive, and neuter. A transitive verb is a verb that expresses an action which has some person or thing for its object: as, An intransitive verb is a verb that expresses an action which has no person or thing for its object: as, A passive verb is a verb in a passive sentence (passive voice) that represents its subject, or what the nominative expresses, as being acted on: as, In a passive sentence, the action is performed on the subject. These sentences have the same denotative meaning, but their connotative meaning is quite different; active verbs are much more powerful and personal. A neuter verb or impersonal passive verb is a verb that expresses neither action nor passion, but simply being, or a state of being: as, Voice. Voice of speech can be active or passive. Principally in passive voice the same tenses can be used as in active voice. There are two forms of passive voice (the second form is preferred): There are however some things to note. Here active and passive do not really have the same meaning. If for example you describe a picture where people build a house, the first sentence is perfectly correct. The second sentence however will be interpreted as the static perfect of the sentence This is, the house is now ready and not under construction. So the correct passive form is Passive voice can be built quite formally by adhering to some rules. You will however not find normally all tenses as in active voice. Formal rules will lead you to monstrosities like the following, you will certainly never hear (already the active sentence is quite monstrous): Moods. Moods are different forms of the verb, each of which expresses the being, action, or passion, in some particular manner. There are five moods; the infinitive, the indicative, the potential, the subjunctive, and the imperative. The infinitive mood is that form of the verb, which expresses the being, action, or passion, in an unlimited manner, and without person or number: as, The indicative mood is that form of the verb, which simply indicates or declares a thing: as, or asks a question: as, The potential mood is that form of the verb which expresses the power, liberty, possibility, or necessity, of the being, action, or passion: as, The subjunctive mood is that form of the verb, which represents the being, action, or passion, as conditional, doubtful, and contingent: as, The imperative mood is that form of the verb which is used in commanding, exhorting, entreating, or permitting: as, Conjugation. The conjugation of a verb is a regular arrangement of its moods, tenses, persons, numbers, and participles. An auxiliary, or a sign of a verb, is a short verb prefixed to one of the morphological forms of another verb, to express some particular mode and time of the being, action, or passion. The auxiliaries are "do, be, have, shall, will, may, can", and "must", with their variations. "Do", "be", and "have" express the indicative mood. Most often, the auxiliaries are used in the following way: Shall and will. Often confused with each other in modern English. These auxiliaries have distinct meanings, and, as signs of the future, they are interchanged thus: Present tense, sign of the indicative first-future. Past tense, sign of aorist, or indefinite. See also: Shall and will by Wikipedia Must. If "must" is ever used in the sense of the past tense, the form is the same as that of the present: this word is entirely invariable. Is being. English grammar has changed, no longer means the same as The first sentence refers to an ongoing action, the second to a completed one. Forms of conjugation. Verb may be conjugated in four ways: The verbs would be conjugated affirmatively, unless said otherwise. Love, conjugated in simple form. The verb "love" is a regular active verb. Simple form, active or neuter. The simplest form of an English conjugation, is that which makes the present and past tenses without auxiliaries; but, even in these, auxiliaries are required for the potential mood, and are often preferred for the indicative. Infinite mood. The infinitive mood is that form of the verb, which expresses the being, action, or passion, in an unlimited manner, and without person or number. It is used only in the present and perfect tenses. Present tense. This tense is the root, or radical verb; and is usually preceded by the preposition to, which shows its relation to some other word: thus, Perfect tense. This tense prefixes the auxiliary have to the past participle; and, like the infinitive present, is usually preceded by the preposition to: thus, Indicative mood. The indicative mood is that form of the verb, which simply indicates or declares a thing, or asks a question. It is used in all the tenses. Present tense. The present indicative, in its simple form, is essentially the same as the present infinitive, or radical verb; except that the verb "be" has "am" in the indicative. The simple form of the present tense is varied thus: This tense may also be formed by prefixing the auxiliary "do" to the verb: thus, Past tense. This tense, in its simple form is the past; which, in all regular verbs, adds "d" or "ed" to the present, but in others is formed variously. The simple form of the past tense is varied thus: This tense may also be formed by prefixing the auxiliary "did" to the present: thus, Perfect tense. This tense prefixes the auxiliary "have" to the past participle: thus, Past perfect tense. This tense prefixes the auxiliary "had" to the past participle: thus, First-future tense. This tense prefixes the auxiliary "shall or will" to the present: thus, Second-future tense. This tense prefixes the auxiliaries "shall have or will have" to the past participle: thus, Potential mood. The potential mood is that form of the verb, which expresses the power, liberty, possibility, or necessity of the being, action, or passion. It is used in the first four tenses; but the potential past is properly an aorist: its time is very indeterminate: as, Present tense. This tense prefixes the auxiliary "may, can, or must", to the radical verb: thus, Past tense. This tense prefixes the auxiliary "might, could, would, or should", to the radical verb: thus, Perfect tense. This tense prefixes the auxiliaries, "may have, can have, or must have", to the past participle: thus, Past perfect tense. This tense prefixes the auxiliaries, "might have, could have, would have, or should have", to the past participle: thus, Subjunctive mood. The subjunctive mood is that form of the verb, which represents the being, action, or passion, as conditional, doubtful, or contingent. This mood is generally preceded by a conjunction: as, "if, that, though, lest, unless, except". But sometimes, especially in poetry, it is formed by a mere placing of the verb before the nominative: as, It does not vary its termination at all, in the different persons. It is used in the present, and sometimes in the past tense; rarely, and perhaps never properly, in any other. As this mood can be used only in a dependent clause, the time implied in its tenses is always relative, and generally indefinite: as, Present tense. This tense is generally used to express some condition on which a future action or event is affirmed. It is therefore erroneously considered by some grammarians, as an elliptical form of the future. In this tense, the auxiliary "do" is sometimes employed: as, Past tense. This tense, like the past of the potential mood, with which it is frequently connected, is properly an aorist, or indefinite tense; for it may refer to time past, present, or future: as, Imperative mood. The imperative mood is that form of the verb, which is used in commanding, exhorting, entreating, or permitting. It is commonly used only in the second person of the present tense. See, conjugated in simple form. The verb "see" is an irregular active verb. Participles.  Present Past Past Perfect  Seeing. Seen. Having seen. Be, conjugated in simple form. The verb "be" is an irregular neuter verb. Morphological forms.  Present Past Present Participle Past Participle.  Be. Was. Being. Been. Participles.  Present Past Past Perfect  Being. Been. Having been. Read, conjugated in progressive form. The verb "read" is an irregular active verb. Compound or progressive form. Active and neuter verbs may also be conjugated, by adding the present participle to the auxiliary verb "be", through all its changes: as, This form of the verb denotes a continuance of the action or state of being, and is, on many occasions, preferable to the simple form of the verb. Morphological forms of the simple verb.  Present Past Present Participle Past Participle  Read. Read. Reading. Read. Participles.  Present Past Past Perfect  Being reading. ———————— Having been reading. Be loved, conjugated in simple form. The verb "be loved" is a regular passive verb. Form of passive verbs. Passive verbs, in English, are always of a progressive form; being made from transitive verbs, by adding the past participle to the auxiliary verb "be", through all its changes: thus from the active transitive verb "love", is formed the passive verb "be loved". Love, conjugated negatively. Form of negation. A verb is conjugated negatively, by placing the adverb "not" and participles take the negative first: as, not to love, not to have loved; not loving, not loved, not having loved. Love, conjugated interrogatively. Form of question. A verb is conjugated interrogatively, in the indicative and potential moods, by placing the nominative after it, or after the first auxiliary: as, Love, conjugated interrogatively and negatively. Form of question with negation. A verb is conjugated interrogatively and negatively, in the indicative and potential moods, by placing the nominative and the adverb "not" after the verb, or after the first auxiliary: as, Irregular verbs. An irregular verb is a verb that does not form the past and the past participle by assuming "d" or "ed": as, "see, saw, seeing, seen". Of this class of verbs there are about one hundred and ten, beside their several derivatives and compounds. Methods of learning irregular verbs: List of the top irregular verbs: Redundant verbs. A redundant verb is a verb that forms the past or the past participle in two or more ways, and so as to be both regular and irregular: as, "thrive, thrived or throve, thriving, thrived or thriven". Of this class of verbs, there are about ninety-five, beside sundry derivatives and compounds. List of the redundant verbs: Defective verbs. A defective verb is a verb that forms no participles, and is used in but few of the moods and tenses: as, "beware, ought, quoth". List of the defective verbs: A short syntax. The finite verb must agree with its subject, as "The birds fly", except the following cases: the conjunction "and", as "Rhetoric and logic are allied," one person or thing, as "Flesh and blood has not revealed it," empathy, as "Consanguinity, and not affinity, is the ground," "each", "every", or "no", as "No one is the same," and the conjunction "or", as "Fear or jealousy affects him." 

&lt; They are the sum of two numbers before (also known as the Fibonacci-sequence): 2 3 5 8 13 21 34 55 89 That is, 2 + 3 = 5, 3 + 5 = 8, 5 + 8 = 13 and so on... 

These numbers are sorted by reverse alphabet: two, three, six, seven, one, nine, four, five, eight 2 3 6 7 1 9 4 5 8 

This textbook is designed for use with first- and second-year college level physics for engineers and scientists. While the content is not mathematically complicated or very advanced, the students are expected to be familiar with differential calculus and some integral calculus. Unlike the Modern Physics textbook, this textbook will stay with the traditional order in presentation of topics in mechanics, thermodynamics, electromagnetism, and geometric optics. These consist the first two semesters and perhaps first few weeks of the third semester. The topics in modern physics, which can be covered during the third semester in the remaining time, can be presented or read in any order. In keeping with maintaining the orthodox order, we will also maintain the traditional chapter-section organization. A few suggested break points between semesters are shown below as well. These break points are marked by a horizontal line between chapters. "(Note to editors: For the purpose of hierarchical organization, at least until the organization of the book is settled down, it should be: "Physics with Calculus/General Topic/Specific Topic", where "General Topic" is any of the following: Mechanics, Thermodynamics, Waves, Electromagnetism, Optics, Modern. No additional "general topic" should be necessary. "Specific Topic" is what it sounds like it is. It can be as specific as necessary, such as "Conservation of Angular Momentum in Spin-Orbit Coupling", or as general as necessary, such as "Newton's Laws". More specific information, such as ordering of chapters will be kept in this module only. This should minimize the need to rename books each time one section moves from one chapter to another, without the unclear "Part I" or "Unit I".)" 

Very few programs follow exactly one control path and have each instruction stated explicitly. In order to program effectively, it is necessary to understand how one can alter the steps taken by a program due to user input or other conditions, how some steps can be executed many times with few lines of code, and how programs can appear to demonstrate a rudimentary grasp of logic. C constructs known as conditionals and loops grant this power. From this point forward, it is necessary to understand what is usually meant by the word "block". A block is a group of code statements that are associated and intended to be executed as a unit. In C, the beginning of a block of code is denoted with { (left curly), and the end of a block is denoted with }. It is not necessary to place a semicolon after the end of a block. Blocks can be empty, as in {}. Blocks can also be nested; i.e. there can be blocks of code within larger blocks. Conditionals. There is likely no meaningful program written in which a computer does not demonstrate basic decision-making skills. It can actually be argued that there is no meaningful human activity in which some sort of decision-making, instinctual or otherwise, does not take place. For example, when driving a car and approaching a traffic light, one does not think, "I will continue driving through the intersection." Rather, one thinks, "I will stop if the light is red, go if the light is green, and if yellow go only if I am traveling at a certain speed a certain distance from the intersection." These kinds of processes can be simulated in C using conditionals. A conditional is a statement that instructs the computer to execute a certain block of code or alter certain data only if a specific condition has been met. The most common conditional is the If-Else statement, with conditional expressions and Switch-Case statements typically used as more shorthanded methods. Before one can understand conditional statements, it is first necessary to understand how C expresses logical relations. C treats logic as being arithmetic. The value 0 (zero) represents false, and all other values represent true. If you chose some particular value to represent true and then compare values against it, sooner or later your code will fail when your assumed value (often 1) turns out to be incorrect. Code written by people uncomfortable with the C language can often be identified by the usage of #define to make a "TRUE" value. Because logic is arithmetic in C, arithmetic operators and logical operators are one and the same. Nevertheless, there are a number of operators that are typically associated with logic: Relational and Equivalence Expressions:. New programmers should take special note of the fact that the "equal to" operator is ==, not =. This is the cause of numerous coding mistakes and is often a difficult-to-find bug, as the expression codice_1 sets codice_2 equal to codice_3 and subsequently evaluates to codice_3; but the expression codice_5, which is usually intended, checks if codice_2 is equal to codice_3. It needs to be pointed out that, if you confuse = with ==, your mistake will often not be brought to your attention by the compiler. A statement such as codice_8 is considered perfectly valid by the language, but will always assign 20 to codice_9 and evaluate as true. A simple technique to avoid this kind of bug (in many, not all cases) is to put the constant first. This will cause the compiler to issue an error, if == got misspelled with =. Note that C does not have a dedicated boolean type as many other languages do. 0 means false and anything else true. So the following are equivalent:  if (foo()) {  // do something and  if (foo() != 0) {  // do something Often codice_10 and codice_11 are used to work around the lack of a boolean type. This is bad practice, since it makes assumptions that do not hold. It is a better idea to indicate what you are actually expecting as a result from a function call, as there are many different ways of indicating error conditions, depending on the situation.  if (strstr("foo", bar) &gt;= 0) {  // bar contains "foo" Here, codice_12 returns the index where the substring foo is found and -1 if it was not found. Note that this would fail with the codice_13 definition mentioned in the previous paragraph. It would also not produce the expected results if we omitted the codice_14. One other thing to note is that the relational expressions do not evaluate as they would in mathematical texts. That is, an expression codice_15 does not evaluate as you probably think it might. Mathematically, this would test whether or not "value" is between "myMin" and "myMax". But in C, what happens is that "value" is first compared with "myMin". This produces either a 0 or a 1. It is this value that is compared against myMax. Example:  int value = 20;  if (0 &lt; value &lt; 10) { // don't do this! it always evaluates to "true"!  /* do some stuff */ Because "value" is greater than 0, the first comparison produces a value of 1. Now 1 is compared to be less than 10, which is true, so the statements in the if are executed. This probably is not what the programmer expected. The appropriate code would be  int value = 20;  if (0 &lt; value &amp;&amp; value &lt; 10) { // the &amp;&amp; means "and"  /* do some stuff */ Logical Expressions. Here's an example of a larger logical expression. In the statement:  e = ((a &amp;&amp; b) || (c &gt; d)); e is set equal to 1 if a and b are non-zero, or if c is greater than d. In all other cases, e is set to 0. C uses short circuit evaluation of logical expressions. That is to say, once it is able to determine the truth of a logical expression, it does no further evaluation. This is often useful as in the following:  int myArray[12];  if (i &lt; 12 &amp;&amp; myArray[i] &gt; 3) { In the snippet of code, the comparison of i with 12 is done first. If it evaluates to 0 (false), i would be out of bounds as an index to myArray. In this case, the program never attempts to access myArray[i] since the truth of the expression is known to be false. Hence we need not worry here about trying to access an out-of-bounds array element if it is already known that i is greater than or equal to zero. A similar thing happens with expressions involving the or || operator.  while (doThis() || doThat()) ... doThat() is never called if doThis() returns a non-zero (true) value. The If-Else statement. If-Else provides a way to instruct the computer to execute a block of code only if certain conditions have been met. The syntax of an If-Else construct is:  if (/* condition goes here */) {  /* if the condition is non-zero (true), this code will execute */  } else {  /* if the condition is 0 (false), this code will execute */ The first block of code executes if the condition in parentheses directly after the "if" evaluates to non-zero (true); otherwise, the second block executes. The "else" and following block of code are completely optional. If there is no need to execute code if a condition is not true, leave it out. Also, keep in mind that an "if" can directly follow an "else" statement. While this can occasionally be useful, chaining more than two or three if-elses in this fashion is considered bad programming practice. We can get around this with the Switch-Case construct described later. Two other general syntax notes need to be made that you will also see in other control constructs: First, note that there is no semicolon after "if" or "else". There could be, but the block (code enclosed in { and }) takes the place of that. Second, if you only intend to execute one statement as a result of an "if" or "else", curly braces are not needed. However, many programmers believe that inserting curly braces anyway in this case is good coding practice. The following code sets a variable c equal to the greater of two variables a and b, or 0 if a and b are equal.  if (a &gt; b) {  c = a;  } else if (b &gt; a) {  c = b;  } else {  c = 0;  Consider this question: why can't you just forget about "else" and write the code like:  if (a &gt; b) {  c = a;   if (a &lt; b) {  c = b;   if (a == b) {  c = 0; There are several answers to this. Most importantly, if your conditionals are not mutually exclusive, "two" cases could execute instead of only one. If the code was different and the value of a or b changes somehow (e.g.: you reset the lesser of a and b to 0 after the comparison) during one of the blocks? You could end up with multiple "if" statements being invoked, which is not your intent. Also, evaluating "if" conditionals takes processor time. If you use "else" to handle these situations, in the case above assuming (a &gt; b) is non-zero (true), the program is spared the expense of evaluating additional "if" statements. The bottom line is that it is usually best to insert an "else" clause for all cases in which a conditional will not evaluate to non-zero (true). The conditional expression. A conditional expression is a way to set values conditionally in a more shorthand fashion than If-Else. The syntax is:  (/* logical expression goes here */) ? (/* if non-zero (true) */) : (/* if 0 (false) */) The logical expression is evaluated. If it is non-zero (true), the overall conditional expression evaluates to the expression placed between the ? and :, otherwise, it evaluates to the expression after the :. Therefore, the above example (changing its function slightly such that c is set to b when a and b are equal) becomes:  c = (a &gt; b) ? a : b; Conditional expressions can sometimes clarify the intent of the code. Nesting the conditional operator should usually be avoided. It's best to use conditional expressions only when the expressions for a and b are simple. Also, contrary to a common beginner belief, conditional expressions do not make for faster code. As tempting as it is to assume that fewer lines of code result in faster execution times, there is no such correlation. The Switch-Case statement. Say you write a program where the user inputs a number 1-5 (corresponding to student grades, A(represented as 1)-D(4) and F(5)), stores it in a variable grade and the program responds by printing to the screen the associated letter grade. If you implemented this using If-Else, your code would look something like this:  if (grade == 1) {  printf("A\n");  } else if (grade == 2) {  printf("B\n");  } else if /* etc. etc. */ Having a long chain of if-else-if-else-if-else can be a pain, both for the programmer and anyone reading the code. Fortunately, there's a solution: the Switch-Case construct, of which the basic syntax is:  switch (/* integer or enum goes here */) {  case /* potential value of the aforementioned int or enum */:  /* code */  case /* a different potential value */:  /* different code */  /* insert additional cases as needed */  default:  /* more code */ The Switch-Case construct takes a variable, usually an int or an enum, placed after "switch", and compares it to the value following the "case" keyword. If the variable is equal to the value specified after "case", the construct "activates", or begins executing the code after the case statement. Once the construct has "activated", there will be no further evaluation of "case"s. Switch-Case is syntactically "weird" in that no braces are required for code associated with a "case". Very important: Typically, the last statement for each case is a break statement. This causes program execution to jump to the statement following the closing bracket of the switch statement, which is what one would normally want to happen. However if the break statement is omitted, program execution continues with the first line of the next case, if any. This is called a "fall-through". When a programmer desires this action, a comment should be placed at the end of the block of statements indicating the desire to fall through. Otherwise another programmer maintaining the code could consider the omission of the 'break' to be an error, and inadvertently 'correct' the problem. Here's an example:  switch (someVariable) {  case 1:  printf("This code handles case 1\n");  break;  case 2:  printf("This prints when someVariable is 2, along with...\n");  /* FALL THROUGH */  case 3:  printf("This prints when someVariable is either 2 or 3.\n" );  break; If a "default" case is specified, the associated statements are executed if none of the other cases match. A "default" case is optional. Here's a switch statement that corresponds to the sequence of if - else if statements above. Back to our example above. Here's what it would look like as Switch-Case:  switch (grade) {  case 1:  printf("A\n");  break;  case 2:  printf("B\n");  break;  case 3:  printf("C\n");  break;  case 4:  printf("D\n");  break;  default:  printf("F\n");  break; A set of statements to execute can be grouped with more than one value of the variable as in the following example. (the fall-through comment is not necessary here because the intended behavior is obvious)  switch (something) {  case 2:  case 3:  case 4:  /* some statements to execute for 2, 3 or 4 */  break;  case 1:  default:  /* some statements to execute for 1 or other than 2,3,and 4 */  break; Switch-Case constructs are particularly useful when used in conjunction with user defined "enum" data types. Some compilers are capable of warning about an unhandled enum value, which may be helpful for avoiding bugs. Loops. Often in computer programming, it is necessary to perform a certain action a certain number of times or until a certain condition is met. It is impractical and tedious to simply type a certain statement or group of statements a large number of times, not to mention that this approach is too inflexible and unintuitive to be counted on to stop when a certain event has happened. As a real-world analogy, someone asks a dishwasher at a restaurant what he did all night. He will respond, "I washed dishes all night long." He is not likely to respond, "I washed a dish, then washed a dish, then washed a dish, then...". The constructs that enable computers to perform certain repetitive tasks are called loops. While loops. A while loop is the most basic type of loop. It will run as long as the condition is non-zero (true). For example, if you try the following, the program will appear to lock up and you will have to manually close the program down. A situation where the conditions for exiting the loop will never become true is called an infinite loop.  int a = 1;  while (42) {  a = a * 2; Here is another example of a while loop. It prints out all the powers of two less than 100.  int a = 1;  while (a &lt; 100) {  printf("a is %d \n", a);  a = a * 2; The flow of all loops can also be controlled by break and continue statements. A break statement will immediately exit the enclosing loop. A continue statement will skip the remainder of the block and start at the controlling conditional statement again. For example:  int a = 1;  while (42) { // loops until the break statement in the loop is executed  printf("a is %d ", a);  a = a * 2;  if (a &gt; 100) {  break;  } else if (a == 64) {  continue; // Immediately restarts at while, skips next step  printf("a is not 64\n"); In this example, the computer prints the value of a as usual, and prints a notice that a is not 64 (unless it was skipped by the continue statement). Similar to If above, braces for the block of code associated with a While loop can be omitted if the code consists of only one statement, for example:  int a = 1;  while (a &lt; 100)  a = a * 2; This will merely increase a until a is not less than 100. When the computer reaches the end of the while loop, it always goes back to the while statement at the top of the loop, where it re-evaluates the controlling condition. If that condition is "true" at that instant -- even if it was temporarily 0 for a few statements inside the loop -- then the computer begins executing the statements inside the loop again; otherwise the computer exits the loop. The computer does not "continuously check" the controlling condition of a while loop during the execution of that loop. It only "peeks" at the controlling condition each time it reaches the codice_16 at the top of the loop. It is very important to note, once the controlling condition of a While loop becomes 0 (false), the loop will not terminate until the block of code is finished and it is time to reevaluate the conditional. If you need to terminate a While loop immediately upon reaching a certain condition, consider using break. A common idiom is to write:  int i = 5;  while (i--) {  printf("java and c# can't do this\n"); This executes the code in the while loop 5 times, with i having values that range from 4 down to 0 (inside the loop). Conveniently, these are the values needed to access every item of an array containing 5 elements. For loops. For loops generally look something like this:  for ("initialization"; "test"; "increment") {  /* code */ The "initialization" statement is executed exactly once - before the first evaluation of the "test" condition. Typically, it is used to assign an initial value to some variable, although this is not strictly necessary. The "initialization" statement can also be used to declare and initialize variables used in the loop. The "test" expression is evaluated each time before the code in the "for" loop executes. If this expression evaluates as 0 (false) when it is checked (i.e. if the expression is not true), the loop is not (re)entered and execution continues normally at the code immediately following the FOR-loop. If the expression is non-zero (true), the code within the braces of the loop is executed. After each iteration of the loop, the "increment" statement is executed. This often is used to increment the loop index for the loop, the variable initialized in the initialization expression and tested in the test expression. Following this statement execution, control returns to the top of the loop, where the "test" action occurs. If a "continue" statement is executed within the "for" loop, the increment statement would be the next one executed. Each of these parts of the for statement is optional and may be omitted. Because of the free-form nature of the for statement, some fairly fancy things can be done with it. Often a for loop is used to loop through items in an array, processing each item at a time.  int myArray[12];  int ix;  for (ix = 0; ix &lt; 12; ix++) {  myArray[ix] = 5 * ix + 3; The above for loop initializes each of the 12 elements of myArray. The loop index can start from any value. In the following case it starts from 1.  for (ix = 1; ix &lt;= 10; ix++) {  printf("%d ", ix); which will print  1 2 3 4 5 6 7 8 9 10 You will most often use loop indexes that start from 0, since arrays are indexed at zero, but you will sometimes use other values to initialize a loop index as well. The "increment" action can do other things, such as "decrement". So this kind of loop is common:  for (i = 5; i &gt; 0; i--) {  printf("%d ", i); which yields  5 4 3 2 1 Here's an example where the test condition is simply a variable. If the variable has a value of 0 or NULL, the loop exits, otherwise the statements in the body of the loop are executed.  for (t = list_head; t; t = NextItem(t)) {  /* body of loop */ A WHILE loop can be used to do the same thing as a FOR loop, however a FOR loop is a more condensed way to perform a set number of repetitions since all of the necessary information is in a one line statement. A FOR loop can also be given no conditions, for example:  for (;;) {  /* block of statements */ This is called an infinite loop since it will loop forever unless there is a break statement within the statements of the for loop. The empty test condition effectively evaluates as true. It is also common to use the comma operator in for loops to execute multiple statements.  int i, j, n = 10;  for (i = 0, j = 0; i &lt;= n; i++, j += 2) {  printf("i = %d , j = %d \n", i, j); Special care should be taken when designing or refactoring the conditional part, especially whether using &lt; or &lt;= , whether start and stop should be corrected by 1, and in case of prefix- and postfix-notations. ( On a 100 yards promenade with a tree every 10 yards there are 11 trees. )  int i, n = 10;  for (i = 0; i &lt; n; i++)  printf("%d ", i); // processed n times =&gt; 0 1 2 3 ... (n-1)  printf("\n");  for (i = 0; i &lt;= n; i++)  printf("%d ", i); // processed (n+1) times =&gt; 0 1 2 3 ... n  printf("\n");  for (i = n; i--;)  printf("%d ", i); // processed n times =&gt; (n-1) ...3 2 1 0  printf("\n");  for (i = n; --i;)  printf("%d ", i); // processed (n-1) times =&gt; (n-1) ...4 3 2 1  printf("\n"); Do-While loops. A DO-WHILE loop is a post-check while loop, which means that it checks the condition after each run. As a result, even if the condition is zero (false), it will run at least once. It follows the form of:  do {  /* do stuff */  } while (condition); Note the terminating semicolon. This is required for correct syntax. Since this is also a type of while loop, break and continue statements within the loop function accordingly. A continue statement causes a jump to the test of the condition and a "break" statement exits the loop. It is worth noting that Do-While and While are functionally almost identical, with one important difference: Do-While loops are always guaranteed to execute at least once, but While loops will not execute at all if their condition is 0 (false) on the first evaluation. One last thing: goto. goto is a very simple and traditional control mechanism. It is a statement used to immediately and unconditionally jump to another line of code. To use goto, you must place a label at a point in your program. A label consists of a name followed by a colon (:) on a line by itself. Then, you can type "goto "label";" at the desired point in your program. The code will then continue executing beginning with "label". This looks like:  MyLabel:  /* some code */  goto MyLabel; The ability to transfer the flow of control enabled by gotos is so powerful that, in addition to the simple if, all other control constructs can be written using gotos instead. Here, we can let "S" and "T" be any arbitrary statements:  if ("cond") {  S;  } else {  T;  /* ... */ The same statement could be accomplished using two gotos and two labels:  if ("cond") goto Label1;  T;  goto Label2;  Label1:  S;  Label2: Here, the first goto is conditional on the value of "cond". The second goto is unconditional. We can perform the same translation on a loop:  while ("cond1") {  S;  if ("cond2")  break;  T;  /* ... */ Which can be written as:  Start:  if (!"cond1") goto End;  S;  if ("cond2") goto End;  T;  goto Start;  End: As these cases demonstrate, often the structure of what your program is doing can usually be expressed without using gotos. Undisciplined use of gotos can create unreadable, unmaintainable code when more idiomatic alternatives (such as if-elses, or for loops) can better express your structure. Theoretically, the goto construct does not ever "have" to be used, but there are cases when it can increase readability, avoid code duplication, or make control variables unnecessary. You should consider first mastering the idiomatic solutions, and use goto only when necessary. Keep in mind that many, if not most, C style guidelines "strictly forbid" use of goto, with the only common exceptions being the following examples. One use of goto is to break out of a deeply nested loop. Since break will not work (it can only escape one loop), goto can be used to jump completely outside the loop. Breaking outside of deeply nested loops without the use of the goto is always possible, but often involves the creation and testing of extra variables that may make the resulting code far less readable than it would be with goto. The use of goto makes it easy to undo actions in an orderly fashion, typically to avoid failing to free memory that had been allocated. Another accepted use is the creation of a state machine. This is a fairly advanced topic though, and not commonly needed. Examples. int main(void)  int years;  printf("Enter your age in years : ");  fflush(stdout);  errno = 0;  if (scanf("%d", &amp;years) != 1 || errno)  return EXIT_FAILURE;  printf("Your age in days is %d\n", years * 365);  return 0; 

In C programming, all executable code resides within a function. Note that other programming languages may distinguish between a "function", "subroutine", "subprogram", "procedure", or "method" -- in C, these are all functions. Functions are a fundamental feature of any high level programming language and make it possible to tackle large, complicated tasks by breaking tasks into smaller, more manageable pieces of code. At a lower level, a function is nothing more than a memory address where the instructions associated with a function reside in your computer's memory. In the source code, this memory address is usually given a descriptive name which programmers can use to call the function and execute the instructions that begin at the function's starting address. The instructions associated with a function are frequently referred to as a block of code. After the function's instructions finish executing, the function can return a value and code execution will resume with the instruction that immediately follows the initial call to the function. If this doesn't make immediate sense to you, don't worry. Understanding what is happening inside your computer at the lowest levels can be confusing at first, but will eventually become very intuitive as you develop your C programming skills. For now, it's enough to know that a function and its associated block of code is often executed (called) several times, from several different places, during a single execution of a program. As a basic example, suppose you are writing a program that calculates the distance of a given (x,y) point to the x-axis and to the y-axis. You will need to compute the absolute value of the whole numbers x and y. We could write it like this (assuming we don't have a predefined function for absolute value in any library): /*this function computes the absolute value of a whole number.*/ int abs(int x)  if (x&gt;=0) return x;  else return -x; /*this program calls the abs() function defined above twice.*/ int main()  int x, y;  printf("Type the coordinates of a point in 2-plane, say P = (x,y). First x=");  scanf("%d", &amp;x);  printf("Second y=");  scanf("%d", &amp;y);  printf("The distance of the P point to the x-axis is %d. \n Its distance to the y-axis is %d. \n", abs(y), abs(x));  return 0; The next example illustrates the usage of a function as a procedure. It's a simplistic program that asks students for their grade for three different courses and tells them if they passed a course. Here, we created a function, called codice_1 that can be called as many times as we need to. The function saves us from having to write the same set of instructions for each class the student has taken. /*the 'check' function is defined here.*/ void check(int x)  if (x&lt;60)  printf("Sorry! You will need to try this course again.\n");  else  printf("Enjoy your vacation! You have passed.\n"); /*the program starts here at the main() function, which calls the check() function three times.*/ int main()  int a, b, c;  printf("Type your grade in Mathematics (whole number). \n");  scanf("%d", &amp;a);  check(a);  printf("Type your grade in Science (whole number). \n");  scanf("%d", &amp;b);  check(b);  printf("Type your grade in Programming (whole number). \n");  scanf("%d", &amp;c);  check(c);  /* this program should be replaced by something more meaningful.*/  return 0; Notice that in the program above, there is no outcome value for the 'check' function. It only executes a procedure. This is precisely what functions are for. More on functions. It's useful to conceptualize a function like a machine in a factory. On the input side of the machine, you dump in the "raw materials," or the input data, that you want the machine to process. Then the machine goes to work and and spits out a finished product, the "return value," to the output side of the machine which you can collect and use for other purposes. In C, you must tell the machine exactly what raw materials it is expected to process and what kind of finished product you want the machine to return to you. If you supply the machine with different raw materials than it expects, or if you try to return a product that's different than what you told the machine to produce, the C compiler will throw an error. Note that a function isn't required to take any inputs. It doesn't have to return anything back to us, either. If we modify the example above to ask the user for their grade inside the codice_2 function, there would be no need to pass the grade value into the function. And notice that the codice_2 doesn't pass a value back. The function just prints out a message to the screen. You should be familiar with some basic terminology related to functions: Writing functions in C. It's always good to learn by example. Let's write a function that will return the square of a number. int square(int x)  int square_of_x;  square_of_x = x * x;  return square_of_x; To understand how to write such a function like this, it may help to look at what this function does as a whole. It takes in an int, x, and squares it, storing it in the variable square_of_x. Now this value is returned. The first int at the beginning of the function declaration is the type of data that the function returns. In this case when we square an integer we get an integer, and we are returning this integer, and so we write int as the return type. Next is the name of the function. It is good practice to use meaningful and descriptive names for functions you may write. It may help to name the function after what it is written to do. In this case we name the function "square", because that's what it does - it squares a number. Next is the function's first and only argument, an int, which will be referred to in the function as x. This is the function's "input". In between the braces is the actual guts of the function. It declares an integer variable called square_of_x that will be used to hold the value of the square of x. Note that the variable square_of_x can only be used within this function, and not outside. We'll learn more about this sort of thing later, and we will see that this property is very useful. We then assign x multiplied by x, or x squared, to the variable square_of_x, which is what this function is all about. Following this is a return statement. We want to return the value of the square of x, so we must say that this function returns the contents of the variable square_of_x. Our brace to close, and we have finished the declaration. Written in a more concise manner, this code performs exactly the same function as the above: int square(int x)  return x * x; Note this should look familiar - you have been writing functions already, in fact - main is a function that is always written. In general. In general, if we want to declare a function, we write  "type" "name"("type1" "arg1", "type2" "arg2", ...)  /* "code" */ We've previously said that a function can take no arguments, or can return nothing, or both. What do we write if we want the function to return nothing? We use C's void keyword. void basically means "nothing" - so if we want to write a function that returns nothing, for example, we write void sayhello(int number_of_times)  int i;  for(i=1; i &lt;= number_of_times; i++) {  printf("Hello!\n"); Notice that there is no return statement in the function above. Since there's none, we write void as the return type. (Actually, one can use the return keyword in a procedure to return to the caller before the end of the procedure, but one cannot return a value as if it were a function.) What about a function that takes no arguments? If we want to do this, we can write for example float calculate_number(void)  float to_return=1;  int i;  for(i=0; i &lt; 100; i++) {  to_return += 1;  to_return = 1/to_return;  return to_return; Notice this function doesn't take any inputs, but merely returns a number calculated by this function. Naturally, you can combine both void return and void in arguments together to get a valid function, also. Recursion. Here's a simple function that does an infinite loop. It prints a line and calls itself, which again prints a line and calls itself again, and this continues until the stack overflows and the program crashes. A function calling itself is called recursion, and normally you will have a conditional that would stop the recursion after a small, finite number of steps.  // don't run this! void infinite_recursion()  printf("Infinite loop!\n");  infinite_recursion(); A simple check can be done like this. Note that ++depth is used so the increment will take place before the value is passed into the function. Alternatively you can increment on a separate line before the recursion call. If you say print_me(3,0); the function will print the line Recursion 3 times. void print_me(int j, int depth)  if(depth &lt; j) {  printf("Recursion! depth = %d j = %d\n",depth,j); //j keeps its value  print_me(j, ++depth); Recursion is most often used for jobs such as directory tree scans, seeking for the end of a linked list, parsing a tree structure in a database and factorising numbers (and finding primes) among other things. Static functions. If a function is to be called only from within the file in which it is declared, it is appropriate to declare it as a static function. When a function is declared static, the compiler will know to compile an object file in a way that prevents the function from being called from code in other files. Example: static int compare( int a, int b )  return (a+4 &lt; b)? a : b; Using C functions. We can now "write" functions, but how do we use them? When we write main, we place the function outside the braces that encompass main. When we want to use that function, say, using our calculate_number function above, we can write something like  float f;  f = calculate_number(); If a function takes in arguments, we can write something like  int square_of_10;  square_of_10 = square(10); If a function doesn't return anything, we can just say  say_hello(); since we don't need a variable to catch its return value. Functions from the C Standard Library. While the C language doesn't itself contain functions, it is usually linked with the C Standard Library. To use this library, you need to add an #include directive at the top of the C file, which may be one of the following from C89/C90: The functions available are: Variable-length argument lists. Functions with variable-length argument lists are functions that can take a varying number of arguments. An example in the C standard library is the printf function, which can take any number of arguments depending on how the programmer wants to use it. C programmers rarely find the need to write new functions with variable-length arguments. If they want to pass a bunch of things to a function, they typically define a structure to hold all those things -- perhaps a linked list, or an array -- and call that function with the data in the arguments. However, you may occasionally find the need to write a new function that supports a variable-length argument list. To create a function that can accept a variable-length argument list, you must first include the standard library header stdarg.h. Next, declare the function as you would normally. Next, add as the last argument an ellipsis ("..."). This indicates to the compiler that a variable list of arguments is to follow. For example, the following function declaration is for a function that returns the average of a list of numbers:  float average (int n_args, ...); Note that because of the way variable-length arguments work, we must somehow, in the arguments, specify the number of elements in the variable-length part of the arguments. In the average function here, it's done through an argument called n_args. In the printf function, it's done with the format codes that you specify in that first string in the arguments you provide. Now that the function has been declared as using variable-length arguments, we must next write the code that does the actual work in the function. To access the numbers stored in the variable-length argument list for our average function, we must first declare a variable for the list itself:  va_list myList; The va_list type is a type declared in the stdarg.h header that basically allows you to keep track of your list. To start actually using myList, however, we must first assign it a value. After all, simply declaring it by itself wouldn't do anything. To do this, we must call va_start, which is actually a macro defined in stdarg.h. In the arguments to va_start, you must provide the va_list variable you plan on using, as well as the name of the last variable appearing before the ellipsis in your function declaration: float average (int n_args, ...)  va_list myList;  va_start (myList, n_args);  va_end (myList); Now that myList has been prepped for usage, we can finally start accessing the variables stored in it. To do so, use the va_arg macro, which pops off the next argument on the list. In the arguments to va_arg, provide the va_list variable you're using, as well as the primitive data type (e.g. int, char) that the variable you're accessing should be: float average (int n_args, ...)  va_list myList;  va_start (myList, n_args);  int myNumber = va_arg (myList, int);  va_end (myList); By popping n_args integers off of the variable-length argument list, we can manage to find the average of the numbers: float average (int n_args, ...)  va_list myList;  va_start (myList, n_args);  int numbersAdded = 0;  int sum = 0;  while (numbersAdded &lt; n_args) {  int number = va_arg (myList, int); // Get next number from list  sum += number;  numbersAdded += 1;  va_end (myList);  float avg = (float)(sum) / (float)(numbersAdded); // Find the average  return avg; By calling average (2, 10, 20), we get the average of 10 and 20, which is 15. 

Having covered the basic concepts of C programming, we can now briefly discuss the process of "compilation". Like any programming language, C by itself is completely incomprehensible to a microprocessor. Its purpose is to provide an intuitive way for humans to provide instructions that can be easily converted into machine code that "is" comprehensible to a microprocessor. The compiler is what translates our human-readable source code into machine code. To those new to programming, this seems fairly simple. A naive compiler might read in every source file, translate everything into machine code, and write out an executable. That could work, but has two serious problems. First, for a large project, the computer may not have enough memory to read all of the source code at once. Second, if you make a change to a single source file, you would have to recompile the "entire" application. To deal with these problems, compilers break the job into steps. For each source file (each codice_1 file), the compiler reads the file, reads the files it references via the codice_2 directive, and translates them to machine code. The result of this is an "object file" (codice_3). After all the object files are created, a "linker" program collects all of the object files and writes the actual executable program. That way, if you change one source file, only that file needs to be recompiled, after which, the application will need to be re-linked. Without going into details, it can be beneficial to have a superficial understanding of the compilation process. Preprocessor. The preprocessor provides the ability for the inclusion of so called header files, macro expansions, conditional compilation and line control. These features can be accessed by inserting the appropriate preprocessor directives into your code. Before compiling the source code, a special program, called the preprocessor, scans the source code for tokens, or special strings, and replaces them with other strings or code according to specific rules. The C preprocessor is not technically part of the C language and is instead a tool offered by your compiler's software. All preprocessor directives begin with the hash character (#). You can see one preprocessor directive in the Hello world program. Example:  #include &lt;stdio.h&gt; This directive causes the stdio header to be included into your program. Other directives such as codice_4 control compiler settings and macros. The result of the preprocessing stage is a text string. You can think of the preprocessor as a non-interactive text editor that modifies your code to prepare it for compilation. The language of preprocessor directives is agnostic to the grammar of C, so the C preprocessor can also be used independently to process other kinds of text files. Syntax Checking. This step ensures that the code is valid and will sequence into an executable program. Under most compilers, you may get messages or warnings indicating potential issues with your program (such as a conditional statement always being true or false, etc.) When an error is detected in the program, the compiler will normally report the file name and line that is preventing compilation. Object Code. The compiler produces a machine code equivalent of the source code that can be linked into the final program. At this point the code itself can't be executed, as it requires linking to do so. It's important to note after discussing the basics that compilation is a "one way street". That is, compiling a C source file into machine code is easy, but "decompiling" (turning machine code into the C source that creates it) is not. Decompilers for C do exist, but the code they create is hard to understand and only useful for reverse engineering. Linking. Linking combines the separate object files into one complete program by integrating libraries and the code and producing either an or a . Linking is performed by a linker program, which is often part of a compiler suite. Common errors during this stage are either missing or duplicate functions. Automation. For large C projects, many programmers choose to automate compilation, both in order to reduce user interaction requirements and to speed up the process by recompiling only modified files. Most Integrated Development Environments (IDE's) have some kind of project management which makes such automation very easy. However, the project management files are often usable only by users of the same integrated development environment, so anyone desiring to modify the project would need to use the same IDE. On UNIX-like systems, make and Makefiles are often used to accomplish the same. Make is traditional and flexible and is available as one of the standard developer tools on most Unix and GNU distributions. Makefiles have been extended by the GNU Autotools, composed of Automake and Autoconf for making software compilable, testable, translatable and configurable on many types of machines. Automake and Autoconf are described in detail in their respective manuals. The Autotools are often perceived to be complicated and various simpler build systems have been developed. Many components of the GNOME project now use the declarative Meson build system which is less flexible, but instead focuses on providing the features most commonly needed from a build system in a simple way. Other popular build systems for programs written in the C language include CMake and Waf. Once gcc is installed, it can be called with a list of c source files that have been written but not yet compiled. e.g. if the file main.c includes functions described in myfun.h and implemented in myfun_a.c and myfun_b.c, then it is enough to write  gcc main.c myfun_a.c myfun_b.c myfun.h is included in main.c, but if it is in a separate header file directory, then that directory can be listed after a "-I " switch. In larger programs, Makefiles and gnu make program can compile c files into intermediate files ending with suffix .o which can be linked by gcc. How to compile each object file is usually described in the Makefile with the object file as a label ending with a colon followed by two spaces (tabs often cause problems) followed by a list of other files that are dependencies, e.g. .c files and .o files compiled in another section, and on the next line, the invocation of gcc that is required. Typing codice_5 or codice_6 often gives the information needed to on how to use make, as well as gcc. Although gcc has a lot of option switches, one often used is -g to generate debugging information for gdb to allow gdb to show source code during a step-through of the machine code program. gdb has an 'h' command that shows what it can do, and is usually started with 'gdb a.out' if a.out is the anonymous executable machine code file that was compiled by gcc. 

Introduction. Recall from the previous section that we considered the case where F["x"]/&lt;m("x")&gt; analogous to modular arithmetic but with polynomials, and that when we are looking at numbers modulo "n", we have a field iff Zn is a field if "n" is prime. Can we say something similar about F["x"]/&lt;m("x")&gt;? Indeed, if m("x") is irreducible then F["x"]/&lt;m("x")&gt; is a field. This section deals with these kinds of fields, known as a finite field. Definitions. We have the object F["x"]/&lt;m("x")&gt; where this is the set of polynomials in F["x"] are divided by the polynomial m("x"). Of the elements in F["x"]/&lt;m("x")&gt; we can easily define addition, subtraction, multiplication, division and so on normally but with a reduction modulo m("x") to get the desired remainder. We have that F["x"]/&lt;m("x")&gt; is a commutative ring with identity, and if m("x") is irreducible then F["x"]/&lt;m("x")&gt; is a field. If m("x") has degree "n", then If F is Zp (so "p" is prime) then formula_2 Properties. Now remember with complex numbers C, we have "invented" the symbol i to stand for the root of the solution "x"2+1=0. In fact, we have C=R["x"]/&lt;"x"2+1&gt;. When we have a "general" finite field, we can do this also. We write this often as F["x"]/&lt;m("x")&gt;=F(α) where α is "the root of" m("x") - we "define" α to be the root of m("x"). F(α) in fact is the smallest field which contains F and α. Finite field theorems. We have a number of theorems associated with finite fields. Some examples. Let's look at a few examples that go through these ideas. The complex numbers. Complex numbers, briefly, are numbers in the form where "i" is the solution to the equation "x"2+1=0 These numbers in fact form a field, however it is not a finite field. Take m("x")="x"2+1, with the field F being R. Then we can form the complex numbers as F/&lt;m("x")&gt;. Now F/&lt;m("x")&gt; = { "a"+"bx" | "a", "b" ∈ R} because the remainders must be of degree less than m("x") - which is 2. So then ("a"+"bx")("c"+"dx")="ac"+"bdx"2+("ad"+"bc")"x". But remember that we are working in F/&lt;m("x")&gt;. So "x"2 modulo "x"2+1, can be written as ("x"2+1)-1=-1, and substituting -1 above yields a rather familiar expression. If we let the symbol "i" to be the "root of "x"2+1", then "i"2+1=0 and "i"2=-1. The rest of the field axioms follow from here. We can then say the complex numbers C=R/&lt;"x"2+1&gt;=R("i"). The Zp case. We can still do this for some field in general. Let's take Z3 for example, and pick m("x")="x"2+"x"+2. m("x") is irreducible - m(0)=2, m(1)=4=1, m(2)=4+2+2=8=2. So Z3/&lt;"x"2+"x"+2&gt; is a finite field. Assume α is a root of m("x"). Then Z3(α) = { "a"+"b"α|"a", "b" ∈ Z3}. Since Z3/&lt;"x"2+"x"+2&gt; is finite, we can list out all its elements. We have the constant terms, then the α terms, then the α+constant terms, and so on. We have {0, 1, 2, α, α+1, α+2, 2α, 2α+1, 2α+2}. Now we have α2+α+2=0, then We can verify multiplication works mod m("x") - for example Reducing coefficients normally mod 3 we get Now using the formula above for α2 Verify for yourself that multiplication and other operations work too. Primitive elements. Recall from ../Modular arithmetic/ that the order of a number "a" modulo "n", in a field, is the least power such that "a""k"-1=1 in Zn, with "k" the size of this field. Since the order is defined over a field, this leads us to consider whether we have primitive elements in F["x"]/&lt;m("x")&gt; - which we do. If we have F(α), just like in Zn, α is primitive iff the order of α is "q"-1 where "q" is the number of elements in F["x"]/&lt;m("x")&gt;. Let's take Z2/&lt;"x"2+"x"+1&gt;. Is α (root of "x"2+"x"+1) primitive? First, if α is a root of "x"2+"x"+1, Now, let us calculate powers of α Recall that the size of this field is 4 (the "n" in Zn, in this case, 2, raised to the power of the degree of the polynomial, in this case 2). Now we have α3=α4-1=1, and α is primitive. We generally want to look at powers of α in F(α), to see whether they are primitive, since we already know about the orders of the constants in F(α) - which we've looked at in ../Modular arithmetic/. 

The word biology means, "the science of life", from the Greek bios, "life", and logos, "word" or "knowledge." Therefore, Biology is the science of Living Things. That is why Biology is sometimes known as Life Science. The science has been divided into many subdisciplines, such as botany, bacteriology, anatomy, zoology, histology, mycology, embryology, parasitology, genetics, molecular biology, systematics, immunology, microbiology, physiology, cell biology, cytology, ecology, and virology. Other branches of science include or are comprised in part of biology studies, including paleontology, taxonomy, evolution, phycology, helimentology, protozoology, entomology, biochemistry, biophysics, biomathematics, bio engineering, bio climatology and anthropology. Characteristics of life. Not all scientists agree on the definition of just what makes up life. Various characteristics describe most living things. However, with most of the characteristics listed below we can think of one or more examples that would seem to break the rule, with something nonliving being classified as living or something living classified as nonliving. Therefore we are careful not to be too dogmatic in our attempt to explain which things are living or nonliving. An easy way to remember this is GRIMNERD C All organisms; - Grow, Respire, Interact, Move, Need Nutrients, Excrete (Waste), Reproduce,Death, Cells (Made of) Living things are organized in the microscopic level from atoms up to cells. Atoms are arranged into molecules, then into macromolecules, which make up organelles, which work together to form cells. Beyond this, cells are organized in higher levels to form entire multicellular organisms. Cells together form tissues, which make up organs, which are part of organ systems, which work together to form an entire organism. Of course, beyond this, organisms form populations which make up parts of an ecosystem. All of the Earth's ecosystems together form the diverse environment that is the earth. Example:- sub atoms, atoms, molecules, cells, tissues, organs, organ systems, organisms, population, community, eco systems Emergent property is viewed in the biological organization of life, ranging from the subatomic level to the entire biosphere. Emergent properties are not unique to life, but biological systems are far more complex, making the emergent properties of life difficult to study. Systems biology is a biology-based inter-disciplinary field of study that focuses on complex interactions within biological systems, using a holistic approach. Biologists study properties of life, with reductionist approach and holistic approach. Nature of science. Science is a methodology for learning about the world. It involves the application of knowledge. The scientific method deals with systematic investigation, reproducible results, the formation and testing of hypotheses, and reasoning. Reasoning can be broken down into two categories, induction (specific data is used to develop a generalized observation or conclusion) and deduction (general information leads to specific conclusion). Most reasoning in science is done through induction. Science as we now know it arose as a discipline in the 17th century. Scientific method. The scientific method is not a step by step, linear process. It is an intuitive process, a methodology for learning about the world through the application of knowledge. Scientists must be able to have an "imaginative preconception" of what the truth is. Scientists will often observe and then hypothesize the reason why a phenomenon occurred. They use all of their knowledge and a bit of imagination, all in an attempt to uncover something that might be true. A typical scientific investigation might go like so: Scientists first make observations that raise a particular question. In order to explain the observed phenomenon, they develop a number of possible explanations, or hypotheses. This is the inductive part of science, observing and constructing plausible arguments for why an event occurred. Experiments are then used to eliminate one or more of the possible hypotheses until one hypothesis remains. Using deduction, scientists use the principles of their hypothesis to make predictions, and then test to make sure that their predictions are confirmed. After many trials (repeatability) and all predictions have been confirmed, the hypothesis then may become a theory. Quick Definitions The scientific method is based primarily on the testing of hypotheses by experimentation. This involves a control, or subject that does not undergo the process in question. A scientist will also seek to limit variables to one or another very small number, single or minimum number of variables. The procedure is to form a hypothesis or prediction about what you believe or expect to see and then do everything you can to violate that, or falsify the hypotheses. Although this may seem unintuitive, the process serves to establish more firmly what is and what is not true. A founding principle in science is a lack of absolute truth: the accepted explanation is the most likely and is the basis for further hypotheses as well as for falsification. All knowledge has its relative uncertainty. Theories are hypotheses which have withstood repeated attempts at falsification. Common theories include evolution by natural selection and the idea that all organisms consist of cells. The scientific community asserts that much more evidence supports these two ideas than contradicts them. Charles Darwin. Charles Darwin is most remembered today for his contribution of the theory of evolution through natural selection. The seeds of this theory were planted in Darwin's mind through observations made on a five-year voyage through the New World on a ship called the Beagle. There, he studied fossils and the geological record, geographic distribution of organisms, the uniqueness and relatedness of island life forms, and the affinity of island forms to mainland forms. Upon his return to England, Darwin pondered over his observations and concluded that evolution must occur through natural selection. He declined, however, to publish his work because of its controversial nature. However, when another scientist, Wallace, reached similar conclusions, Darwin was convinced to publish his observations in 1859. His hypothesis revolutionized biology and has yet to be falsified by empirical data collected by mainstream scientists. After Darwin. Since Darwin's day, scientists have amassed a more complete fossil record, including microorganisms and chemical fossils. These fossils have supported and added subtleties to Darwin's theories. However, the age of the Earth is now held to be much older than Darwin thought. Researchers have also uncovered some of the preliminary mysteries of the mechanism of heredity as carried out through genetics and DNA, areas unknown to Darwin. Another growing area is comparative anatomy including homology and analogy. Today we can see a bit of evolutionary history in the development of embryos, as certain (although not all) aspects of development recapitulate evolutionary history. The molecular biology study of slowly mutating genes reveal considerable evolutionary history consistent with fossil and anatomical record. Challenges to Darwin. Darwin and his theories have been challenged many times in the last 150 years. The challenges have been primarily religious based on a perceived conflict with the preconceived notion of creationism. Many of those who challenge Darwin have been adherents to the young earth hypothesis that says that the Earth is only some 6000 years old and that all species were individually created by a god. Some of the proponents of these theories have suggested that chemical and physical laws that exist today were different or nonexistent in earlier ages. However, for the most part, these theories are either not scientifically testable and fall outside the area of attention of the field of biology, or have been disproved by one or more fields of science. References. "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of Cleveland State University." 

Puzzles | How do you... ? | Ten Apples and A Basket Ten Apples and A Basket puzzle You have a basket containing ten apples. You have ten friends, who each desire an apple. You give each of your friends one apple. After ten minutes, each of your friends has one apple each. Yet there is one remaining apple in the basket. How is this possible? Solution 

&lt; Back to Problem Solutions: //Solutions below are only humorous and not Logical // 

Matter. Matter is defined as anything that has mass (an amount of matter in an object) and occupies space (which is measured as volume). Water. Hydrogen bonding. Water organizes nonpolar molecules Ionization of water: H2O -&gt; H+ + OH- pH Buffer References. "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of the Cleveland State University." 

Building blocks of life Hereditary (Genetic) information. RNA DNA origin. Proteins: Their building block is amino acids. The bond connecting 2 of the amino acids together are called peptide bonds. One of these bonds makes a monopeptide, two a dipeptide, and any more than that makes a polypeptide. References. "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of the Cleveland State University." 

Origin of life: 3 hypotheses. -However, this would imply that some other origin of life was likely because it would have had to happen elsewhere before it could be transported here, and the only difference would be that life did not originate on Earth. -This is often criticized for the improbability of life being produced by a chemical reaction caused by lightning, along with the ability of any produced DNA to be in a sequence that could reproduce a working life model. It is also attacked for religious reasons, as it bypasses things like the idea of a supreme being directly creating humans. It also seems unlikely to some that such huge changes are possible in evolution without evidence of a credible "in-between” stage. Many of the stages of man are disputed due to their somewhat shaky grounds. For example, bones from other animals or species have been mistaken for humanoid while complete skeletons have been put together from a limited number of bones. -This is often attacked for many of the same reasons that religion is attacked, and is often regarded as superstitious and/or unscientific. The early earth. It is believed that the Earth was formed about 4.5 billion years ago. Chemical reactions on early earth Issues Origin of cells. Cells are very small and decompose quickly after death. As such, fossils of the earliest cells do not exist. Scientists have had to form a variety of theories on how cells (and hence life) was created on Earth. The earliest cells. Cyanobacteria References. "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of the Cleveland State University." 

Biological membranes. Biological membranes surround cells and serve to keep the insides separated from the outsides. They are formed of phospholipid bilayers, which by definition are a double layer of fatty acid molecules (mostly phospholipids, lipids containing lots of phosphorus). Proteins serve very important functions in cellular membranes. They are active transports in and out of the cell, acting as gatekeepers. They relay signals in and out of the cell. Proteins are the site of many enzymatic reactions in the cell, and play a role in regulation of cellular processes. Phospholipid. Phospholipid bilayer Membrane proteins. Integral membrane proteins Channel protein Transport across membranes Osmotic pressure Bulk transport Receptor-mediated endocytosis. Active transport Cotransport (coupled transport) "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of the Cleveland State University." 

=Cell-cell interactions= Cells interact with the environment and with each other. Communicating junctions. Gap junctions. "This text is based on notes very generously donated by Dr. Paul Doerder, Ph.D., of the Cleveland State University." 

Interlingua is an international auxiliary language that uses simple grammar and vocabulary common among Romance and Germanic languages. This means that it is a great option for speakers of those languages looking to learn another, or for speakers of other languages who are looking for a nice gateway in English, Italian, Spanish, French, German, etc. 



Biochemical pathways. Evolution of biochemical pathways References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Glucose + O2 → CO2 + H2O + ATP Glycolysis overview. Glycolysis accounting 2NADH  4 ATP (from 2 G3P) –2 ATP (priming)  2 ATP (net gain) Summary: The net input of glycolysis is 2 ATP molecules which are used to split one glucose molecule. The net yield of this step is 2 ATP and 2 pyruvate. Lactate formation. Either lactic acid or alcohol can be formed as a result of anaerobic respiration in cells. 

 6 CO2 + 6 H2O → C6H12O6 + 6 O2 Light Reactions. Accessory pigments. Photosynthetic steps The Even More Detailed Light Reactions. What the Light Reactions Do: The light reactions of photosynthesis occur in chloroplasts in and on the thylakoid disks. During the light reactions, light energy charges up ATP molecules. More specifically, light turns the chloroplast into an acid battery, and this battery charges up ATP. How the "Chloroplast-Battery" Charges ATP: The stroma is the fluid inside of the chloroplasts, and it carries a negative charge. This means that it contains about a "gazillion" extra electrons. The solvent of stroma is water. The fluid inside the thylakoid disks is positively charged because it contains a lot of hydrogen (H+) ions. The pH here is low, making the fluid very acidic. The solvent of thylakoid disk fluid is water. A chloroplast acts like a battery, because it has separated a strong positive charge and a strong negative charge in two different compartments. Energy is released when H+ ions (free protons) flow from the inside of a thylakoid disk to the stroma. This is electrical energy, since it is a flow of charged particles. The protons pass through special channels (made of protein) in the thylakoid membrane; this reaction is 'exothermic.' The energy that is given off is used to fuel this reaction (Pi is the phosphate ion):  ADP + Pi --&gt; ATP The proton can go to the negative stroma, but only if it uses its energy to charge up ATP. Since one reaction wants to go, and the other one doesn't, and since the first reaction releases energy and the second one absorbs energy, the two reactions are known to be 'coupled' together so that the first fuels the second. Of course, a special enzyme must be involved for this to happen. Chlorophyll Molecules on a Thylakoid Disk: Hundreds of chlorophyll molecules cover the surface of a thylakoid disk, making the disk green. The nonpolar "tails" of the chlorophyll molecule are embedded in the membrane of the thylakoid. “Dark” reactions. The Detailed Dark Reactions. What the Dark Reactions Do: The dark reactions build sugar from carbon dioxide gas (CO2), water (H2O), and energy from ATP molecules that were charged up during the light reactions. The dark reactions occur in the stroma of a chloroplast. Dark reactions usually occur in the light, but they don't have to. They'll occur in the dark until the chloroplast's supply of ATP runs out (usually about 30 seconds). The Calvin Cycle: The Calvin Cycle is the fancy name for the metabolic pathway that builds sugar. This means that it involves a whole lot of chemical reactions, and it uses a lot of different enzymes to catalyze the reactions. Carbon dioxide gas is stable, therefore the bonds that hold the carbon and oxygen atoms are strong. Therefore it takes a lot of energy to break the bonds and separate the carbon atoms from the oxygen atoms. The energy needed to do this comes from ATP molecules. When inorganic carbon (like from CO2) is being added to an organic molecule (such as sugar), this is called carbon fixation. It takes 2 complete turns of the Calvin Cycle to make a glucose molecule. References. "Some portions of this text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." The detailed portions are not provided by Dr. Doerder. 

How cells divide Bacterial DNA replication. Eukaryotic chromosomes Chromosome number. Chromosome numbers Human chromosomes. Cell cycle Mitotic cell cycle.  - at this point the chromosomes are composed of two sister chromotids connected by a common centromere. Mitosis. Kinetochore Microtubules attach to kinetochores. Metaphase Anaphase Plant mitosis. Cell cycle control Molecular control of cell cycle: Cdk and cyclin References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Sexual life cycle. Typical animal life cycle References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Charles Darwin, for all he contributed to the science of biology, never knew about the mechanism by which living things inherit traits from previous generations, or how new traits arise. As any schoolchild can tell you, this mechanism of interitance has since been found to be DNA, or deoxyribonucleic acid. DNA allows for stable inheritance of traits: the code in each strand of DNA is replicated precisely through the pairing of basic units along each strand. The error rate in this replication is amazingly low; not even one base pair in a million matches out of sequence. However, when even one base pair is added to a new strain of DNA in an order differently than in the parent chain, it can be the basis of a mutation. These changes in DNA sequences are the microscopic origin of changes in traits of all studied living things. Even the smallest difference in a strand of DNA can result in a change in traits that can cost the life of the organism. Mutations can produce proteins with a new or altered function. In humans, the example of Sickle cell anemia is commonly given as its origin is a difference of only one base pair in a section of DNA that encodes red blood cells. Individual sequences of DNA that encode for specific proteins are called genes and are the units of heredity. Each one has a set nucleotide, and together all of the genes (and some sequence of DNA that apparently do not code for any biologically important functions) together make up the entire chromosome Modern Y chromosome. Y-chromosome is the most evolved chromosome. Generally it is thought that if Y- chromosome is present in an individual then he will be male. But if mutation occurs at sex determining region or zinc factor then it will not code for testis determining factor, and results in normal female. This type of female's frequency is 1/250000. X-chromosome inactivation. In females, one X-chromosone is randomly switched off forming a Barr body. Barr body. Dense region in the nucleus formed by the inactive X-chromosome. Human genetic disorders. Down's Syndrome (Mongolism) Down's Syndrome is usually produced by the nondisjunction of chromosome 21 during oogenesis and sometimes during spermatogenesis. The individual suffering from this type of syndrome has 47 chromosomes instead of the normal 46. The extra chromosome is not a sex chromosome but an autosome. Most cases of mongolism were found to occur in children born by women in their forties. The affected children, called mongoloids, show mental retardation and have a shorter life expectancy. Their most prominent feature is the Mongolian folds in their eyes; hence, the term mongolism. Klinefelter's Syndrome When an XY-bearing sperm unites with an X-bearing egg, the resulting condition is called Klinefelter's Syndrome, or sexually undeveloped male. Individuals having the syndrome show the following characteristics: The same abnormal meiotic division may occur in females. They produce eggs with XX or no sex chromosomes. Such egg, when fertilized by a Y-bearing sperm, will not develop (YO). This is because YO is lethal—it wil cause death to the offspring. 

DNA. DNA stands for "Deoxyribose Nucleic Acid". That is, a nucleic acid with two sugars. DNA is the hereditary material of cells and is considered the blueprint of life. DNA is found in all kingdoms of life. Even most viruses have DNA. A molecule of DNA is chemically stable (it does not have a 2-prime alcohol group.) When someone says DNA, they may be referring to one's genetic material on multiple levels: They may be speaking about a single deoxyribose nucleic acid molecule, a section of a double helix, a section of a chromosome, or one's entire hereditary composition. Hershey-Chase Experiment. The Hershey and Chase experiment was one of the leading suggestions that DNA was a genetic material. Hershey and Chase used phages, or viruses, to implant their own DNA into a bacterium. They did two experiments marking either the DNA in the phage with a radioactive phosphorus or the protein of the phage with radioactive sulfur. With the bacteria that was infected by the phages with radioactive DNA the DNA in the bacteria was radioactive. In the bacteria that was infected with the radioactive protein the bacteria was radioactive, not the DNA. This proves that DNA is a genetic material and it is passed on in viruses. DNA/RNA components. Structure of DNA. DNA is in a double helix structure made up of nucleotides. The "backbone" of the double helix is composed of phosphates connected to a five-carbon sugar called "deoxyribose". The "rungs" are composed of organic compounds, purines and pyrimidines. Purines contain Adenine(A) and Guanine(G) and have two rings in their structures. Pyrimidines contain Cytosine(C) and Thymine (T) and have one ring in their structures. Franklin. DNA model Replicon. A region of DNA that is replicated from a single origin. References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Flow of genetic information Transcription bubble. Promoter site Translation in bacteria. tRNA References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Pembentukan Kata "(Word Forming)". In Indonesian, you can form a new word from an existing word. For example, in English, you form the word "organize" from the word "organ". Further, you can add "-ization" suffix on it and thus "organization". This example pretty much illustrates how Indonesian words are formed. Certainly, this formation is not necessarily the same and it's far richer than that. You'll learn shortly why. The word formation is very attached to Indonesian / Malay culture. Some of the formations may not make sense to the westerners. However, with some stretch of imagination, you should be able to cope with it. Many Asian language speakers, on the other hand, will find it very intuitive and much more structured than, say, Chinese. Using these formations, you just need to memorize one type of word (either the noun form, the verb form, or the adjective form). You can then "nounize" the verb or "verbize" the noun or form other interesting words. In a sense, the word formation is the irregular side of Indonesian since you must know which word works with which affixes. This can be confusing for beginners. However, even if you use the incorrect affixes, as long as it doesn't lead to preexisting meaning, sympathetic speakers would understand. Imbuhan spesifik "(Specific Affixes)". I call these very specific affixes because, unlike the previous ones, they can only be used in very few words. 

A form of Indonesian Conjugation Awalan me- "(Prefix me-)". Prefix "me-" is used to form active verbs. You may combine any nouns, any adjective, numbers, and even verbs themselves to form another verb. Prefixing any words with "me-" would require a minor spelling change (inflection) in order to facilitate a smooth transition in pronunciation. The inflection is based solely on the first letter of the original word. This is an example of consonant mutation. See the table below. Note: Me- + Kata Kerja "(Me- + Verb)". Verbs can be combined with "me-" prefix. The meaning is exactly the same as the infinitive. The semantics may be different depending on the type of verbs. To refresh our mind -- there are two kinds of verbs in Indonesian: Me- + Kata Kerja Transitif "(Me- + Transitive Verbs)". Transitive verbs can not be used in a sentence in their infinitive forms, except for a few words (e.g. "makan" (= to eat), "minum" (= to drink)). So, in order to use it, you must conjugate it with "me-" or some other prefixes. Example: Me- + Kata Kerja Intransitif "(Me- + Intransitive Verbs)". Unlike transitive verbs, only a few of intransitive verbs can be conjugated with "me-". Example: Most intransitive verbs must be used in their infinitive forms (e.g. "tidur" (= to sleep)) or use other affixes. Note that, the sense of transitivity is not the same with other languages like Spanish. Here the word "tari" (= to dance) is considered intransitive. Me- + Kata Benda "(Me- + Noun)". When "me-" prefix is combined with a noun, the new word could mean one of the following: Note: The trickiest part of "me-" is when it is being combined with noun. Not all nouns can be combined with "me-" and the meaning is highly dependent on the culture. Me- + Kata Sifat "(Me- + Adjective)". When combined with adjective, the "me-" prefix indicates the subject changes more into or to turn to the indicated adjective. For example: Note: Adjective that describes emotions cannot be combined with "me-". Me- + Kata Bilangan "(Me- + Numbers)". The only numbers that can be conjugated with "me-" are "satu" (= one) and "dua" (= two). Me- + Kata Tempat "(Me- + Places)". When combined with places, "me-" would indicate that the subject is going to the said place. Me- + Kata Seru "(Me- + Expletive)". Note that the "expletives" doesn't mean "swear words" here, but rather the words that mimic sounds. For example: "Moo", "Baa", "Quack", etc. When expletives are combined with me-, it indicates that the subject produces the indicated sound. 



Introduction. You are already familiar with writing a number down, and having it mean a certain combination of tens, hundreds, and so on. This is one form of number representation, but there are others. We will look at number bases and continued fractions. Number bases. You are already familiar with base-10 number representation. For example, the number 2818 is the same as We can replace the number 10 here with any number and we obtain a different number. In general, we can represent an integer "n" in a base "b" by the following: where the "a"i are all less than "b". We write a number base "b" as ("a"k"a"k-1..."a"0)"b". For example, if we have the numeral 243 in base 6, we write it (243)6. When we are in base 10 we simply write the number: for example the numeral 155 in base 10 is simply written 155. However, given a number in a base "b", how can we convert it to our natural base 10 system? Likewise, how can we convert a number from our base 10 system to base "b"? The first is relatively easy, the other more difficult. Converting base "b" to base 10. We simply use the definition of a base-"b" number to convert a base-"b" number to base 10 - that is we simply multiply out. For example Converting base 10 to base "b". This task however is slightly more difficult, and there are several ways of performing this task. One method is to write each step using the division algorithm from the ../Number theory/ section. For example, if we wish to convert 1317 to base 12: So in base 12, (919)12=1317. Real numbers. We've just seen how we can convert "integers" from base to base, but how do we work with converting "real" numbers? Recall in base 10 we write a number such as 11.341 as and so we can extend our definition above of a base-"b" number to be "a"k"b"k+"a"k-1"b"k-1+...+"a"0"b"0+"a"-1"b"-1+... where the "a"i are all less than "b", and the sum could terminate or go on indefinitely. Again, how are we to convert these numbers from base to base? We can convert the integral part, but what about the "fractional part" (the part less than 1)? Converting fractional "n" to base-"b". Say we wish to convert .341 in base 10 to base 6. We write a table, using the following rules  "i" "c"i γi 6γi  0 0 .341 2.046  1 2 .046 0.276  2 0 .276 1.656  3 1 .656 3.936  4 3 .936 5.616  5 5 .616 3.696  6 3 .696 4.176  7 4 .176 1.056  8 1 .056 0.336  9 0 .336 2.016 It looks like this expansion will go on forever! We need to calculate for further values of i to see whether we hit a repeat value of γi to see whether we get a repetition. So we have the approximation that .341 is near to (.20135341)6. ("Calculate this using the definition. How close is our approximation?") If we obtain a base-"b" representation for example, that looks something like (.18191819181918191819...)b we call the representation "periodic". We often write this as We use this same procedure to convert a fractional number to base-"b" by replacing the 6 above with "b". Converting fractional "n" to base 10. We have a nifty trick we can use to convert a fractional "n" in base-"b" to base 10 provided the representation repeats. Let us look at an example: Consider formula_4. Now then And now which is Then And finally On converting (3145)7 to base 10, we obtain the following Continued fractions. In a sense, the base-"b" representation is nice, but it has a few shortcomings in respect to accuracy. We cannot reliably represent the number formula_12 using base-"b" representation. This is where the "continued fraction" representation comes in handy, which has some nice properties regarding quadratic irrationals. A "continued fraction" is a number in the form Since this notation is clearly rather cumbersome, we abbreviate the above to Again we ask ourselves how can we convert from and to this number representation? Again converting from is simpler than converting to. Converting from continued fraction representation. We simply use our definition of the continued fraction to convert from a continued fraction. This may look difficult, but in fact is quite simple depending on which end one starts with. Let's look at an example Consider Now, we work from right to left. We first have the fraction The next digit 2 tells us to perform and then take the reciprocal to get Now the next digit 1 tells us to perform and then to take the reciprocal to get Now we must add "q"0 which is always greater than 1 and we get the result: Converting to continued fraction representation. Again, we draw up a table. Consider the fraction 12/22, and use the following rules in the table  i θi θi-1 qi  0 12/22 . 0  1 12/22 22/12 1  2 5/6 6/5 1  3 1/5 5/1 5 So now the continued fraction for 12/22 is [0; 1, 1, 5]. Converting a periodic continued fraction to quadratic irrational. Firstly, we make note of a nice property of periodic continued fractions (where the sequence of "q"i repeat), that For example, consider the continued fraction Now which we can rewrite as Now we can simplify to obtain which is a quadratic equation and can be solved to obtain Note that we can always roll up the continued fraction into the form (aα+b)/(cα+d)=α, which demonstrates the point that every quadratic irrational has a repeating continued fraction representation Convergents. Say we have a continued fraction ["q"0; "q"1, ... ] which represents a number "n". Let us examine the following series of fractions ["q"0], ["q"0; "q"1], ["q"0; "q"1, "q"2] and so on. Each element of the series is known as a "convergent". It turns out that the series of convergents provide the best rational approximations to "n". These can be calculated from the continued fraction representation, but also from the calculation table. Let us take formula_39. Continue as before, but place an extra -1 row, and set u-1=1, v-1=0. Iterate with the rules and the sequence repeats and the continued fraction is [2; 2, 4, 2, 4, ... ]. We can continue the process to generate more convergents - the convergents are 2, 5/2, 22/9, 49/20, ... 

DNA grooves. Categories of transcription factors in eukaryotes Lac operon of E. coli. lac repressor Regulation in eukaryotes References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

A mutation is a permanent change to an organism's genetic material (DNA or RNA). Mutations are a rare but significant biological process, since they provide the variation on which evolution acts and are also the source of cancer. An organism's genetic material is made up of polymers (chains) of four different nucleotides, like a recipe book written in a language of only four letters. A mutation event is when the order of the nucleotides in DNA change, usually when the DNA is being copied. Mutations come in a number of forms: Point Mutations. Point mutations are all mutations which involve a single nucleotide. These come in the form of substitutions, insertions and deletions: Substitution. Substitution Mutations: In substitution mutations, a nitrogenous base of a triplet codon of DNA is replaced by another nitrogen base or some derivative of the nitrogen base, changing the codon. The altered codon codes for a different amino acid substitution.The substitution mutations are of two types: 1.Transitions: It is the replacement of one purine in a polynucleotide chain by another purine(A by G or G by A) or one pyrimidine by another pyrimidine(T by C or C by T) 2.Transversions:A base pair substitution involving the substitution of a purine by pyrimidine or pyrimidine by a purine is called transversion. Larger mutations. Larger mutations which involve more than one nucleotide also include insertions and deletions, but can also include inversions, rearrangement of nucleotides and duplication of entire genes: Chromosomal mutations. Chromosomal mutations involve changes to entire chromosomes. These mutations are particularly rare: Effects of mutations. Mutations can have a variety of different effects depending on the type of mutation, the significance of the piece of genetic material affected and whether the cells affected are germ-line cells. Only mutations in germ-line cells can be passed on to children, while mutations elsewhere can cause cell-death or cancer. Mutations can be classified by their effects: Silent Mutation. Silent Mutations are DNA mutations that do not result in a change to the amino acid sequence or a protein. They may occur in a non-coding region (outside of a gene or within an intron), or they may occur within an exon in a manner that does not alter the final amino acid chain. Missense Mutation. Missense mutations are types of point mutations where a single nucleotide is changed to cause substitution of a different amino acid. This in turn can render the resulting protein nonfunctional. Such mutations are responsible for diseases such as Epidermolysis bullosa. References. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

=Recombinant DNA technology = Restriction endonucleases. Originally found in bacteria to prevent invasion of viral DNA, cuts double stranded DNA that is unmethylated, will not cut newly synthesized DNA since hemi-methylated, a product of semi-conservative replication of DNA Restriction endonucleases. Gene cloning Automated sequencing. Typical machine Biochips. Microarray chips contain wells of DNA that code for specific genes that uses the concept of hybridization with the gene of interest to see if a gene is expressed or is present. Stem cells. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Fundamental classification Classification of Living Things and Naming of Organisms. In science, the practice of classifying organisms is called taxonomy (Taxis means arrangement and nomos mean method). The modern taxonomic system was developed by the Swedish botanist Carolus Linnaeus (1707-1778). He used simple physical characteristics of organisms to identify and differentiate between different species and is based on genetics. Linnaeus developed a hierarchy of groups for taxonomy. To distinguish different levels of similarity, each classifying group, called taxon (pl. taxa) is subdivided into other groups. To remember the order, it is helpful to use a mnemonic device. The taxa in hierarchical order: The domain is the broadest category, while species is the most specific category available. The taxon Domain was only introduced in 1990 by Carl Woese, as scientists reorganise things based on new discoveries and information. For example, the European Hare would be classified as follows: Eukaryote --&gt; Animalia --&gt; Chordata --&gt; Mammalia --&gt; Lagomorpha --&gt; Leporidae --&gt; Lepus --&gt; "Lepus europaeus". Eukaryote is just one of the classes inside of the main class Phyla Binomial nomenclature. Binomial nomenclature is used to name an organism, where the first word beginning with a capital is the genus of the organism and the second word beginning with lower-case letter is the species of the organism. The name must be in italics and in Latin, which was the major language of arts and sciences in the 18th century. The scientific name can be also abbreviated, where the genus is shortened to only its first letter followed by a period. In our example, "Lepus europaeus" would become "L. europaeus". Taxonomy and binomial nomenclature are both specific methods of classifying an organism. They help to eliminate problems, such as mistaken identity and false assumptions, caused by common names. An example of the former is the fact that a North American robin is quite different from the English robin. An example of the latter is the comparison between crayfish and catfish, where one might believe that they both are fish when in fact, they are quite different. Nomenclature is concerned with the assignment of names to taxonomic groups in agreement with published rules. To study for a test these are the best words to know taxonomist, biologist, chemist, geologist, unicellular, multicellular, bilateral symmetry, radial symmetry, chlorophyll, photosynthesis, respiration, reproduction, vertebrates, endoskeleton, exoskeleton, consumers, decomposers, heterotroph, autotroph, vascular, non-vascular. These are all part of classifying things. Eukaryotes and Prokaryotes. Recall that there are two basic types of cells: eukaryotes and prokaryotes. Eukaryotes are more complex in structure, with nuclei and membrane-bound organelles. Some characteristics of eukaryotes are: Prokaryotes refer to the smallest and simplest type of cells, without a true nucleus and no membrane-bound organelles. Bacteria fall under this category. Some characteristics: Three Domains and Six Kingdoms. The three domains are organized based on the difference between eukaryotes and prokaryotes. Today's living prokaryotes are extremely diverse and different from eukaryotes. This fact has been proven by molecular biological studies (e.g. of RNA structure) with modern technology. The three domains are as follows: Archaea (Archaebacteria) consists of archaebacteria, bacteria which live in extreme environments. The kingdom Archaea belongs to this domain. Eubacteria consists of more typical bacteria found in everyday life. The kingdom Eubacteria belongs to this domain. Eukaryote encompasses most of the world's visible living things. The kingdoms Protista, Fungi, Plantae, and Animalia fall under this category. The Six Kingdoms. Under the three domains are six kingdoms in taxonomy: Animalia, contains general animals and is the largest kingdom with over 1 000 000 species. Plantae, contains all plants on Earth. Protista, the third kingdom, was introduced by the German biologist Ernst Haeckel in 1866 to classify micro-organisms which are neither animals nor plants. Since protists are quite irregular, this kingdom is the least understood and the genetic similarities between organisms in this kingdom are largely unknown. For example, some protists can exhibit properties of both animals and plants. Fungi are organisms that obtain food by absorbing materials in their bodies. Mushrooms and molds belong in this kingdom. Originally, they were part of the plant kingdom but were recategorized when they were discovered not to photosynthesize. Eubacteria are bacteria, made up of small cells, which differ in appearance from the organisms in the above kingdoms. They lack a nucleus and cell organelles. They have cell walls made of peptidoglycan. Archae (or Archaebacteria) are bacteria that live in extreme environments, such as salt lakes or hot, acidic springs. These bacteria are in their own category as detailed studies have shown that they have unique properties and features (ex. unusual lipids that are not found in any other organism)which differ them from other bacteria and which allow them to live where they live. Their cell walls lack peptidoglycan. Origins of Diversity. The diversity in our planet is attributed to diversity within a species. As the world changed in climate and in geography as time passed, the characteristics of species diverged so much that new species were formed. This process, by which new species evolve, was first described by British naturalist Charles Darwin as natural selection. For an organism to change, genetic mutations must occur. At times, genetic mutations are accidental, as in the case of prokaryotes when they undergo asexual reproduction. For most eukaryotes, genetic mutations occur through sexual reproduction, where meiosis produces haploid gametes from the original parent cells. The fusion of these haploid gametes into a diploid zygote results in genetic variation in each generation. Over time, with enough arrangement of genes and traits, new species are produced. Sexual reproduction creates an immense potential of genetic variety. One goal of taxonomy is to determine the evolutionary history of organisms. This can be achieved by comparing species living today with species in the past. The comparison in anatomy and structure is based on data from development, physical anatomy, biochemistry, DNA, behaviour, and ecological preferences. The following are examples of how such data is used: Although a horse and a human may look different, there is evidence that their arm structures are quite similar. Their arms' sizes and proportions may be different, but the anatomical structures are quite similar. Such evidence reveals that animals in different taxa may not be that different. Biological features from a common evolutionary origin are known as homologous. Biochemical analysis of animals similar in appearance have yielded surprising results. For example, although guinea pigs were once considered to be rodents, like mice, biochemistry led them to be in their taxon of their own. Phylogeny, Cladistics, and Cladograms. Modern taxonomy is based on many hypotheses' of the evolutionary history of organisms, known as phylogeny. As with the Scientific Method, scientists develop a hypothesis on the history of an animal and utilise modern science and technology to prove the phylogeny. Cladistics is a classification system which is based on phylogeny. Expanding on phylogeny, cladistics is based on the assumption that each group of related species has one common ancestor and would therefore retain some ancestral characteristics. Moreover, as these related species evolve and diverge from their common ancestor, they would develop unique characteristics. Such characteristics are known as derived characteristics The principles of phylogeny and cladistics can be expressed visually as a cladogram, a branching diagram which acts as a family (phylogenetic) tree for similar species. A cladogram can also be used to test alternative hypotheses for an animal's phylogeny. In order to determine the most likely cladogram, the derived characteristics of similar species are matched and analysed. Classification of Living Things Practice Questions. 1. If taxonomists had to select an existing kingdom to reclassify, which of the six would most likely be chosen? Why? 2. Complete the following without consulting external sources: a) The species "caudatum" is in the family "Paramecidae". What would be the binomial name of this organism? b) Give the abbreviation of the binomial name. 3. a) Irish moss belongs to the genus "Chondrus". The name for this species is "crispus". Give the binomial name. b) Give the abbreviation of the binomial name. 4. Humans and chimpanzees are alike. Which of the following data would most accurately prove this correct? 5. Which of the following is out of order? 6. A taxonomist discovers Organism A and Organism B and wishes to classify them. Which of the following choices is the most informative? 7. DNA analysis is usually done using DNA found in a cell's mitochondria, and not in a cell's nucleus. From your knowledge of mitosis, explain why this is so. 1. Archaebacteria 3.a) Chondrus crispus b) C. cripus 4. B 5. G 6. B 

Introduction. Viruses are the smallest biological particle (the tiniest are only 20 nm in diameter). However, they are not biological organisms so they are not classified in any kingdom of living things. They do not have any organelles and cannot respire or perform metabolic functions. Viruses are merely strands of DNA or RNA surrounded by a protective protein coat called a capsid. Viruses only come to life when they have invaded a cell. Outside of a host cell, viruses are completely inert. Since first being identified in 1935, viruses have been classified into more than 160 major groups. Viruses are classified based on their shape, replication properties, and the diseases that they cause. Furthermore, the shape of a virus is determined by the type and arrangement of proteins in its capsid. Viruses pathogenic to humans are currently classified into 21 groups. Viruses can also attack bacteria and infect bacterial cells. Such viruses are called bacteriophages. Viral Replication. As previously stated, viruses are not a biological life form so they cannot reproduce by themselves. They need to take over a functioning eukaryotic or prokaryotic cell to replicate its DNA or RNA and to make protein coat for new virus particles. In order to enter a cell, a virus must attach to a specific receptor site on the plasma membrane of the host cell. The proteins on the surface of the virus act as keys which fit exactly into a matching glycoprotein on the host cell membrane. In some viruses, the attachment protein is not on the surface of the virus but is in the capsid or in the envelope. There are two forms of viral replication: the lytic cycle and the lysogenic cycle. Lysogenic Cycle. The reproduction cycle of viruses with RNA and no DNA is slightly different. A notable example of a RNA-based virus is HIV, a retrovirus. Viral Genome. The genome of a virus consists of DNA or RNA, whose size and configuration vary. The entire genome can exist as a single nucleic acid molecule or several nucleic acid segments. Also, the DNA or RNA may be single-stranded or double-stranded, and either linear or circular. Not all viruses can reproduce in a host cell by themselves. Since viruses are so small, the size of their genome is limiting. For example, some viruses have coded instructions for only making a few different proteins for the viruses' capsid. On the other hand, the human genome codes for over 30,000 different proteins. Therefore, the lack of coded instructions cause some viruses to need the presence of other viruses to help them reproduce themselves. Such viruses are called replication defective. Lastly, it is worthy to note that 70% of all viruses are RNA viruses. As the process of RNA replication (with enzymes and other organelles of the host cell) is more prone to errors, RNA viruses have much higher mutation rates than do DNA viruses. Viruses Practice Questions. Answers to Viruses Practice Questions 

For Eubacteria, please visit General Biology/Classification of Living Things/Eubacteria. Domains of life: characteristics. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 

Introduction. Out of the six kingdoms, Protista is the most diverse. This is the kingdom of organisms with strange, atypical characteristics. In essence, this kingdom is designated for organisms which do not belong in any other kingdom. The majority of protists are microscopic. Classification of Protists. There are three phyla of protists, based on their type of nutrition. 1. Protozoa (animal-like protists) are heterotrophs that ingest or absorb their food and helps. 2. Algae (plant-like protists) are autotrophs they get nutrition from photosythesis. 3. Slime moulds and water moulds (fungus-like protists) are also heterotrophs, like protozoa. Protozoa. As heterotrophs, protozoa scavenge materials from their surroundings. Others are predators which actively hunt or ambush small organisms such as bacteria and other protozoa for a source of nutrition. Protozoa can be parasitic as well; they may live inside larger organisms, like humans. Most protozoa live as single cells, although a few form colonies. Protozoa are generally difficult to identify due to their varied shape. They may appear as jelly-like blobs, spherical sunbursts, or a flattened leaf. Tiny blood parasites may be only 2 μm long. On the other hand, shell-covered marine may be 5 cm or more in diameter. Furthermore, different protozoans have their own complex life cycles. The complexity has led certain organisms to be mistakenly classified for other species. Nevertheless, protozoa can move, and so, they are classified based on their methods of locomotion. Characteristics of Protozoa : they love environment and each other... Algae. Algae are much simpler than protozoa. They are aquatic and contain chlorophyll. Algae can exist as a single cell or as giant seaweeds 60 m in length. Formerly, algae were classified as plants but this was incorrect as algae lack parts of true plants: leaves, stems, roots, xylem, and phloem (though some plants, such as mosses, also lack all of these parts). Since algae belong in the kingdom Protista, algae is a broad term used to denote all aquatic eukaryotes which photosynthesise; algae can differ in size and shape as well. There are six phyla of algae: chlorophytes (green algae), phaeophytes (brown algae), rhodophytes (red algae), chrysophytes (diatoms), pyrrophytes (dinoflagellates), and euglenophytes (euglenoids). Chlorophytes. Chlorophytes resemble plants the most. Like plants, their cell walls contain cellulose and they store food in reserve as starch. Chlorophytes can be unicellular or multicellular. Most chlorophytes use flagellae for some locomotion. Phaeophytes. Phaeophytes are nearly all multicellular marine organisms, which are known to us as seaweeds. They have cell walls composed of cellulose and alginic acid (a substance similar to pectin). The cellulose and alignic acid help to retain water and prevent seawood from drying out when exposed to air at low tide. Since phaeophytes live in a tidal environment, they have large, flat fronds (a large leaf) which can withstand pounding by waves. Their bases strongly anchor the algae to the rocky seabed and prevent them from being washed out to sea. Phaeophytes are usually found in areas of cold water. Rhodophytes. Rhodophytes are typically found in warmer seawater, and are more delicate and smaller than brown algae (phaeophytes). Rhodophytes are also able to grow at deeper depths in the ocean, since red algae absorb green, violet, and blue light, the wavelengths of which penetrate the deepest below the water surface. They also have mucilaginous material to resist drying. Chryosophytes. Chryosophytes are the most abundant unicellular algae in the oceans. They are also one of the biggest components of plankton, a free-floating collection of microorganisms, eggs, and larvae. As photosynthetic organisms, they produce a significant amount of atmospheric oxygen. The reproduction cycle of chryosophytes is particularly interesting. Note that diatoms reproduce both asexually and sexually. Since diatoms have a rigid cell wall with an outer layer of silica (found in sand and glass), the daughter cells produced by mitosis must fit inside the original cell wall. Therefore, each generation of diatoms is smaller than the one before. The reduction in size continues until the diatoms produce sexually, producing a zygote which eventually grows to the original size as it matures. Pyrrophytes. Pyrrophytes are unicellular, photosynthetic, and mostly aquatic. They have protective coats composed of stiff cellulose. They are more easily identifiable, due to the presence of two flagellae. The longer flagellae propels the dinoflagellate, while the second shorter, flatter flagellae functions as a rudder. Some species of pyrrophytes are zooxanthellae. Since they lack cellulose plates, they make their home in coral reefs and animals, such as sea anemones, and molluscs. In returning the favour of sheltering them, dinoflagellates provide carbohydrates to their host through photosynthesis. This is why there are nutrient-rich coral reefs in malnutritions water. A negative aspect of pyrrophytes is that under certain conditions, species of dinoflagellates reproduce rapidly to form a harmful algal bloom (HAB), known as a red tide if dinoflagellates are the cause. Such pyrrophytes can produce toxins which may injure or kill wildlife, and additionally any consumers of contaminated wildlife. Euglenophytes. Like pyrrophytes, euglenophytes are small unicellular freshwater organisms with two flagella. They are mainly autotrophic or heterotrophic, depending if they have a red, light-sensitive structure called an eyespot. Slime molds &amp; Water molds. There are two phyla of slime moulds and one phylum of water moulds. Oomycotes (Water moulds). Oomycotes are filamentous organisms which resemble fungi, in that they live as saprotrophs. Oomycotes differ from other moulds with the presence of spores and their sexual life cycle. Myxomycetes (Plasmodial slime moulds). Myxomycetes are visible to the naked eye as tiny slug-like organisms which creep over decayed and dead matter. This streaming blob containing many nuclei is called a plasmodium. Acrasiomycetes (Cellular slime moulds) and its reproductive cycle. Acrasiomycetes exist as individual amoeboid cells with one nucleus each. When in unfavourable conditions, each acrasiomycota cell gathers together to form a pseudoplasmodium. Reproductive Cycle:. 1. One acrasiomycota cell joins with others to form a pseudoplasmodium. 2. The pseudoplasmodium shrinks and forms a smaller plasmodium. 3. The plasmodium migrates to a suitable environment. 4. The plasmodium develops a sporangia, where original parental nuclei has divided by meiosis into haploid spores to be germinated. 5. When favourable conditions arise, the spores germinate and are carried away by animals or the wind. 6. Cycle repeats. Protists Practice Questions. 1. Which of the following adjectives describe the major food source of protozoa? 2. The protozoan "Giardia lamblia" can inhabit a human body's intestinal tract and cause gastroenteritis. a) Give the abbreviated binomial name of this protozoan. b) Would the relationship between this protozoan and human being be mutualistic, commensalistic, or parasitic? 3. Found in many products, such as Petri dishes, agar is made from mucilagnious material in seaweed. Of the six phyla of algae, which phyllum/phyla would agar be made from? 4. Which of the following adjectives describe the major food source of Euglenophytes without an eyespot? 5. Can coral reefs exist in nutrient-poor areas? Explain.  

Introduction. Although you may not recognise fungi, they are just as prevalent as plants and animals. Their spores are in the air which we breathe, fungi allow us to make bread, and mushrooms (a type of fungi) are eaten by us. A few types of fungi are unicellular. For example, yeasts live as individual oval or cylindrical cells. However, the majority of fungi are multicellular. Their bodies are composed of hyphae, a network of fine filaments. In a mushroom, the hyphae are densely packed so it is difficult to see the individual structures when a mushroom is eaten. However, a mushroom is only a specialised reproductive part of the whole fungus. The main part of the fungi is underground in a whole web of hyphae, called a mycelium. In the mycelium, each fungal cell is separated from each other by a septum. Each fungal cell may have one or more nuclei and remains connected to the mycelium because the septa are porous, allowing cytoplasm to flow through the hyphae and fungal cell walls, made of a hard material called chitin. Some fungi do not have septa, and they appear to be large, branching, multinucleate cells. Nutrition. Fungi are saprophytes. When they find a source of food (e.g. dead wood, orange peel) , they decompose it and digest it. The enzymes break down larger organic molecules in the substrate into smaller molecules. These smaller molecules diffuse into the fungus, where they are used to allow growth and repair. Fungi which feed on living cells are parasitic. For example, athlete's foot grows on the human foot. These kinds of fungi produce hyphae called haustoria, which can penetrate host cells without immediately killing them. However, they are friendlier species of fungi. Many fungi live symbiotically with plants or animals. For example, most trees have fungi living in close contact with their roots. In this relationship, known as a mycorrhiza, there are many benefits: Fungal Reproduction. Fungi can reproduce in two ways. Firstly, they may asexually produce through fragmentation. This occurs when pieces of hyphae are broken off, which then grow into new mycelia. The second method is by spores. Spores are lightweight structures and windblown designed to be transported over long distances and by many mediums, such as on the bodies of insects and birds. They are additionally light enough to be blown away for hundreds of kilometers. Spores may be asexual and sexual. Their sexual properties can be analysed to classify the four phylla of fungi. Types of Fungi. Zygospore Fungi (Zygomycetes). This phylum includes bread molds and other saprotrophs. Comparable to bacteria, this phylum prefers asexual reproduction over sexual reproduction. 1. Two haploid hyphae of opposite types, also known as mating strain + and mating strain -, combine and fuse together. 2. Plasmogamy, the union of the two parent hyphae, occurs and results in the creation of a heterokaryotic (n + n) zygosporangium or zygospore. Note that the zygospore is NOT diploid yet; the haploid nuclei are simply clumped together. 3. Immediately, a thick wall develops around the zygospore to protect it from drying and other hazards. The zygospore becomes dormant. 4. When conditions are favourable, the zygospore absorbs water and undergoes karyogamy (n + n = 2n), where the haploid nuclei contributed by the two parents fuse to produce diploid zygosporangia. 5. The now diploid zygosporangium then undergoes meiosis to form haploid sporangia. 6. Through asexual reproduction of fungi (See above for more information), the spores from the sporangia germinate and grow into new mycelia. 7. Back to step #1. Club Fungi (Basidiomycetes). This phylum increases mushrooms and shelf fungi. In many ways, the reproduction stages of this phylum is similar to that of zygomycetes. 1. Two haploid hyphae of opposite types, also known as mating strain + and mating strain -, combine and fuse together. 2. Plasmogamy takes place, and a dikaryotic mycelium forms. The dikaryotic mycelium grows faster then the haploid parental mycelia. 3. Environmental factors cause the dikaryotic mycelium to form compact masses which develop into basidiocarps, short-lived reproductive structures. An example is the mushroom. 4. The basidiocarp gills are lined with terminal dikaryotic cells called basidia, which then undergo karyogamy. 5. The basidia are now diploid. They undergo meiosis to develop haploid basidiospores, a term referring to a basidiomycete's spores. 6. Still remaining on the basidiocarp, the haploid basidiospores eject, fall from the basidiocarp, and are dispersed by the wind when mature. 7. In a favourable environment, the basidiospores germinate and grow into short-lived haploid mycelia. 8. Back to Step #1. 

=Multicellular Photosynthetic Autotrophs = Plant phyla. Phyla are 12 groupings Moss life cycle. Moss has no vascular tissues or flowers. It is a thallus plant (it does not have true roots,stem and leaves). It does not produce seeds like the angiosperm. Pterophyta (ferns). Tree fern Fern life cycle Seed plants.  plant Angiosperm life cycle. "This text is based on notes very generously donated by Paul Doerder, Ph.D., of the Cleveland State University." 



Standards. This book is intended to cover science as taught in schools to 14 - 16 year old pupils. It is called GCSE science because the first draft is intended to cover the GCSE (General Certificate of Secondary Education) standard as taught in schools in the UK, but generally in Scotland. Having said that if you are not in the UK and wish to include material to bring it up to another "similar" standard feel free to do so. (Note: It would be appreciated if you could use the "British" way of spelling words (e.g. colour)! Additional resources. There will be teacher notes added in time with worksheets, experiments and additional questions. There will also be a study help desk set up so that students can ask questions. Contributing to this book. If you spot any errors in the texts, just go ahead and edit them out (click the link that says edit this page and start typing). If you wish to help writing any of the modules, either click on any red module and start typing, or if you need help drop me a line on my talk page Authors. Add your name to this section if you have contributed "content" to this book (i.e. not copyediting). Go to contents » 

In networking, a communications protocol or network protocol is the specification of a set of rules for a particular type of communication. Different protocols often describe different aspects of a single communication; taken together, these form a protocol stack. The terms "protocol" and "protocol stack" also refer to the software that implements a protocol. Most recent protocols are assigned by the IETF for internet communications, and the IEEE, or the ISO organizations for other types. The ITU-T handles telecommunications protocols and formats. Index page for network protocols and protocol layers, categorised by the nearest matching layers of the OSI seven layer model. Systems engineering principles have been applied to design network protocols. high Vandyke Common Internet protocols. Common Internet protocols include TCP/IP (Transmission Control Protocol/Internet Protocol), UDP/IP (User Datagram Protocol/Internet Protocol), HTTP (HyperText Transfer Protocol) and FTP (File Transfer Protocol). 

Proxy servers provide a cache of items available on other servers which are presumably slower, more expensive to access or unavailable from the local network. The process of proxying a network through a single host on another network is called network masquerading or IP-masquerading if the source and target networks use the Internet Protocol. This term is used particularly for a World Wide Web server which accepts URLs with a special prefix. When it receives a request for such a URL, it strips off the prefix and looks for the resulting URL in its local cache. If found, it returns the document immediately, otherwise it fetches it from the remote server, saves a copy in the cache and returns it to the requester. The cache will usually have an expiry algorithm which flushes documents according to their age, size, and access history. The Squid cache is the popular http proxy server in UNIX/Linux world. However, Apache's mod_proxy module also provides proxying and caching capabilities, and has the advantage of already being installed on nearly all systems. 

 ^ Polish ^ In Polish it is possible to move words around in the sentence, and to drop subject or object if they are obvious from context. These sentences mean basically the same "Henia has a cat", although they are used in different contexts: Use the first word order, unless emphasizing something. If apparent from the context, you can drop the subject, object or even the verb: In particular, "ja" and "ty", and also their plural equivalents "my" and "wy", are almost always dropped.  ^ Polish ^ 

This book is intended as a compilation of biographies describing the lives and work of influential biologists. 

=Vertebrate digestive system = &lt;br&gt;Functions to break down food into molecules small enough to absorb, or pass across digestive membrane. Digestive tract: tube extending from lips of mouth to anus or cloacae in bird, reptile or monotreme. Lumanal glands: empty into inner body cavity (lumen: inner surface). Tract divided into three main regions: &lt;br&gt;1. buccal cavity &lt;br&gt;2. pharynx &lt;br&gt;3. alimentary canal Alimentary canal divided into four regions: &lt;br&gt;1. esophagus &lt;br&gt;2. stomach &lt;br&gt;3. small intestine &lt;br&gt;4. large intestine Accessory digestive glands, outside digestive tract proper, secrete into lumen of tract through ducts. Includes the salivary glands, liver and pancreas. Buccal cavity, which includes palate and tongue, develops from infolding of stomadeum, or second opening of blastula, whereas the rest of the digestive tract develops from the primitive gut. Teeth: capture and hold prey. In mammals in particular further process and break down food into small particles, increasing surface area available for enzymatic action. Tooth anatomy: &lt;br&gt;1. crown projects above gum, &lt;br&gt;2. root below gum, &lt;br&gt;3. enamel is outer coating of crown, hardest surface in body, of epideral origin &lt;br&gt;4. dentin, below enamel, bone-like and forms bulk of tooth, is harder than bone and contains nerves and blood vessels. (Remember that mammals are heterodontic, with different types of teeth). Pharynx: air passage for adult, gill slits in embryo. Important in lower vertebrates, site of gills. Features derived from pharyngeal region: first pharyngeal pouch gives rise to parts of the ear, other pouches give rise to various other structures. Alimentary canal: epithelium lines lumen, glands secrete into lumen, longitudinal and circular muscles help digestive movements (peristalsis). Esophagus: tube carries food from mouth to stomach. Expands to fit large bolus (lump of chewed food). Secretes mucus for lubrication. Birds have crop for storage, enlargement of esophagus. Epiglottis: keeps food out of air tube, an evolutionary “kludge,” or fix. Stomach. Absorbs water, alcohol, nutrients, uses gastric juice with enzymes, mucous, HCl, released by chief and parietal cells (release protein enzymes) in gastric pits. Rugae: folds of stomach, disappear when full. Sphincter at both ends of stomach, control food passage. Chyme: semi-digested food released to small intestine. Small intestine: three regions, duodenum, jejunum, and ileum. Duodenum site of most intestinal digestion. Jejunum and ileum do most of intestinal absorption. Ileum ends with another sphincter, ileocolic valve or ileosecal valve. Structure: Circular folds covered with villi (singular is villus). Villi: finger-like cellular projections, covered with microvilli, tiny projections which increase surface area. Increases surface area by 900x, speeds digestion (break down) and absorption (taking in nutrients). Large intestine: larger diameter, shorter length than small intestine. No villi. In mammals, forms large gentle loop, colon, empties into straight region, rectum, empties into outside world through anal sphincter. Colon: absorbs water left over, also absorbs vitamins released by bacteria which live there (vitamin K). Food: made up of 1. proteins, 2. fats, 3. carbohydrates 4. fibrous material. Digestive system breaks foods down. Proteins must be broken to amino acids to be absorbed. Polysaccharides to monosaccharides, lipids to fatty acids and monoglycerides to absorb. Salivary glands in mouth, saliva contains mucous, salt and a few enzymes (amalase, begins starch breakdown). Snake venom from oral gland, mixture of toxins and digestive enzymes. Breaks down blood vessels and disables nervous system. Stomach enzymes: released in inactive form, zymogene, converts to active form in lumen of gut. Transformation is triggered by another enzyme, or the stomach’s low pH. Pepsin secreted as pepsinogen (-ogen means primitive form). Stomach glands secrete up to two or three liters a day of gastric juice, which is reabsorbed. Chyme released to duodenum. Small intestine has two major accessory glands: &lt;br&gt;1. pancreas &lt;br&gt;2. liver Pancreas has endocrine and exocrine functions, releases large amounts of carbonate to neutralize acidic chyme, as intestinal enzymes work in neutral pH, and stuff to break down lipids and starch (zymogens, like tripsin) Liver releases bile. Bile made from cholesterol, stored in gall bladder, released in duodenum, emulsifies fats. Emulsify: keeps fats in tiny drops, which are suspended, increasing surface area and action of lipases. Protein and carbohydrates absorbed in intestine, taken to liver for processing. Fatty acids go to lymphatic system Appendix: vestigial remnant. Much variation in digestive systems within mammals: herbivore, carnivore, insectivore, non-ruminant herbivore. Rumen: four-chambered stomach of animals like cows (ruminant herbivores). Cellulose resistant to digestion, rely on microorganisms to break down cellulose. Some bacteria, protists and fungi can break down cellulose, almost no animals can. Bacteria break down cellulose in rumen, to be taken back to the mouth to chew their cud (ruminate). Later cow swallows to proceed with digestion. (Horses not like this). Coprophagy: rabbits and other animals eat their own feces for the nutritious products of the cecum. 

=Circulatory system= Circulatory system functions 1. Transportation&lt;br&gt;  a. Respiration: gas exchange (O2 and CO2), overcomes limited rate of diffusion&lt;br&gt;  b. Nutrition:&lt;br&gt;  c. Excretory: (remove metabolic wastes)&lt;br&gt; 2. Regulation&lt;br&gt;  a. Transport hormones&lt;br&gt;  b. Regulate body temperature&lt;br&gt;  c. Protection&lt;br&gt;  i. Blood clotting&lt;br&gt;  ii. Immune system (carries white blood cells) Vasodilation: allows heat loss across epidermis, as seen in elephant ears, takes more blood to surface of body, sweating may accompany Countercurrent heat exchange: used by dolphins in fins to conserve heat in cold water. Veins surround an artery, and blood returning to body absorbs heat from blood traveling out from body to fin, minimizing heat loss. Used by dogs in feet, etc. Blood made of &lt;br&gt;1. plasma and &lt;br&gt;2. formed cellular elements (red and white blood cells, and platelets). Plasma makes up 55% of blood volume. Cellular elements make up the other 45%. Plasma makeup: 90% water, 7-8% soluble proteins (albumin maintains blood osmotic integrity, others clot, etc.) 1% electrolytes 1% elements in transit Red blood cell (erythrocyte): contains hemoglobin, functions in oxygen transport. In mammals, red blood cells lose nuclei on maturation, and take on biconcave, dimpled, shape. No self repair, live 120 days. About 1000x more red blood cells than white blood cells. About 7-8 micrometers in diameter. Hematocrit: proportion of blood volume that is occupied by cells, about 43% in humans on average. 48% for men and 38% for women. White blood cells (leukocytes): Nucleated, about 10-14 micrometers in diameter, commonly amoeboid, escape circulatory system in capillary beds. Include basophils, eosinophils, neutrophils, monocytes, B- and T-cell lymphocytes. Platelets (thrombocytes) Membrane bound cell fragments in mammals, no nucleus. In non-mammals, platelet role replaced by nucleated cells. Accumulate at site of broken blood vessels, form clots. Bud off special cells in bone marrow. 1-2 micrometers in diameter. 7-8 day life span, 1/10 or 1/20 as abundant as white blood cells. Arteries: carry blood away from heart. Smallest tubes called arterioles, feed blood to capillaries. Veins: return blood to heart. Smallest veins called venules. Structure of arteries and veins, listed from inside (lumen) out: &lt;br&gt;1. epithelium (endothelium), &lt;br&gt;2. elastic connective tissue fibers, &lt;br&gt;3. smooth muscle, &lt;br&gt;4. connective tissue. Arteries have thicker elastic layer than do veins. Capillaries, where exchange of materials occurs, are very thin and narrow, and red blood cells pass through single file. Capillaries are tiny but numerous, and their total volume is greater than that of supplying arteries. Blood velocity drops in capillaries, picks back up in veins. Pressure highest in arteries, lower in capillaries and arteries. Osmotic pressure draws interstitial fluid from blood in arterioles, but replaces it in venules. One-way valves mean that blood can flow only one way, works with residual blood pressure and compression by skeletal muscles. Low pressure in thoracic cavity caused by breathing also helps move blood. 

Lymphatic system: part of the immune system, a one-way, or open, system. Takes up interstitial fluid not taken up by venules. Lymphatic structures:&lt;br&gt; 1. lymphatic capillaries&lt;br&gt; 2. lymphatic vesicles&lt;br&gt; 3. lymph nodes&lt;br&gt; 4. lymphatic organs (spleen and thymus) Lymph: movement in mammals through one-way valves, similar to blood movement in veins. (Some non-mammals have lymphatic hearts of unknown embryonic origin. Frogs and salamanders have several.) Lymph rejoins cardiovascular system into a large vein near the heart via single large thoracic duct. As lymph passes through system, passes lymphocytes, second part of immune system. Heart: pumps blood, design varies between animals. In adult mammal,four chambers form two separate circulations&lt;br&gt; 1. pulmonary circulation to and from lungs and&lt;br&gt; 2. systemic circulation to and from tissues of body. Everything in the heart comes in pairs: 2 atria, 2 ventricles (left and right). Diagrams usually drawn as though animal were on its back. Pattern of blood flow through heart: blood returning from major veins (vena cava) enters right atrium, contraction there delivers blood to right ventricle through a tricuspid valve, one of atrial ventricular valves (AV valve). Contraction of right ventricle drives blood through semi lunar valve into pulmonary circuit and to lungs.Blood return to heart in pulmonary veins, is oxygenated. Goes to left atrium, which contracts and delivers blood to left ventricleby way of aortic semi-lunar valve, then goes to systemic circulation. Both atria and ventricles contract in unison, left is more powerful than right (to all system vs. just lungs). Systole: heart contraction, diastole: heart relaxed Timing of heart contraction: ventricles rebound to relaxed shape (diastole), and semi-lunar valves close. Both atria(singular: atrium) fill with blood coming from pulmonary and systemic circulations.Pressure rises in the atria and blood begins to move into the ventricles.The atria then contract, forcing more blood into the ventricles. There is a pause, then ventricles contract. This raises ventricle pressure, atrio-ventricular(AV) valves shut and semi-lunar valves open, forcing blood from the left ventricle into the major arteries and from the right ventricle into the aorta. Control for this action doesn’t rely on nervous stimulation, has intrinsic rhythmicity, called myogenic. This is the case in mammal as well as in mollusk hearts. Other animals have neurogenic hearts that rely on nervous stimulation for heart action, originating in the cardiac ganglion. The rhythmicity of mammalian heart relies on the sino-atrial (SA)node, or pacemaker. This is a phylogenic (based on evolutionary history) remnant of an early vertebrate heart that had one more chamber than modern hearts. How the heart contracts: waves of depolarization start in SA node and spread through atria. Connectile tissue pauses the spread of depolarization at the atrial ventricular node. Signal continued by bundle branches to lower ventricle, begins to stimulate heart to contract. Contraction starts at bottom of heart at heart apex,then signals spread through heart. Medulla (in the brain) controls autonomic nervous system. (The medulla is part of the brain, is continuous with the spinal cord, and controls involuntary actions of the body). Sympathetic cardiac acceleratorconnects to spinal cord, uses norepinephrine to signal. Parasympathetic cardio-inhibitory center reaches heart through Vagus nerve, usesacetylcholine to signal. Hyperpolarizes membrane to inhibit heart contraction. (Autonomic nervous system: two parts working in contra to control from both sides.) Dominant effect here is inhibitory. If we cut Vagus nerve, heart rate promptly rises about 25 bpm. 

=Respiratory system= In humans and other animals, for example, the anatomical features of the respiratory system include airways, lungs, and the respiratory muscles. Other animals, such as insects, have respiratory systems with very simple anatomical features, and in amphibians even the skin plays a vital role in gas exchange. Plants also have respiratory systems but the directionality of gas exchange can be opposite to that in animals. The respiratory system in plants also includes anatomical features such as holes on the undersides of leaves known as stomata. In mammals, the diaphragm divides the body cavity into the  abdominal cavity: contains the viscera (e.g., stomach and intestines)  thoracic cavity: contains the heart and lungs. Respiratory tree: terminates in alveolus, alveoli. Respiratory bronchioles branch into alveolar ducts and into alveoli. Alveolus: microscopic air sacs, 300 million of these in human lungs. Total surface area large. Gas diffuses micrometer, very tiny distance. 

=Sensory systems= Categorized by Transduction of sensory input into signal. Means to “carry across”,signal transduced, or carried, from environment into nervous signal. Three sensory processes we cover Taste and smell (chemoreception). Found in mammal nose and mouth, fly feet, fish bodies, moth antennae. Papilla: bumps on tongue, contain taste buds down between. Sweet, sour, salty and bitter. Some act directly by ion channels, others act indirectly. Other “taste” sensations really smell. Smell: received in nasopharynx. Airborne molecules go into solution on moist epithelial surface of nasal passage. Approximately 1000 genescode for sensory neuron receptors. “Fried onions” odor not one receptor but a mixture of many odors registered in our mind as one. Very sensitive, habituates rapidly (don’t notice a smell after a bit). Odor sensation has relatively unfiltered root to higher brain centers. Snakes more chemosensory focused than us. Response to gravity and movement. Registered in inner ear. Three semicircular canals loop in three planes at right angles to each other, responsible for transduction of movement messages. Method: hair cells deformed by gelatinous membrane. Vestibular apparatus, gives us perception of gravity and movement. Due to physical response, not chemical binding. Cochlea: bony, coil shaped part of inner ear, where hearing occurs. Sound enters through auditory canal, vibrates tympanic membrane,moving three bones of middle ear (malleus, incus, and stapes)against oval window opening in front of cochlea. Cochlea has three fluid filled ducts, one of these the organ of Corti. Sound waves in air go to vibration in organ of Corti; fluid tickles hair cells, which register the movement along basilar membrane in cochlea. Different sound frequencies move different portions of basilar membrane. Hearing loss due to loss of hair cells.Humans normally smell more than 300 odors in a day(Facts and Truth). Transduction of sound accomplished throgh physical deformation,not chemical binding. Vision. Light enters pupil, focused by lens onto retina. Sclera: hardened part behind retina. Optic nerves and neurons attached to retina. Blind spot where optic nerve attaches, has no receptors. Two types of photoreceptors Fovia: region of most acute vision, has most of the cones, few rods. Transduction process of light to signal a molecular change, to light absorbing molecule called photopigment. Located in outer parts of rods and cones in pigment discs. The rod photopigment is called rhodopsin,cone has three photopigments, called photopsins. This molecular change initiates pathways to result in action potential in downstream neuron leading to vision center in brain.Parul Godika Each of the three photopsins has a different peak of sensitivity: blue,green or red, and changes isometric form (from cisto trans) based on light from a particular wavelength range. Color blindness:inherited lack of one or more types of these cones. Gene carried on X chromosome, therefore more common in men than women. 

Homeostasis. Is a very important part of everyone's and everything's lives. Defined as dynamic constancy of internal environment, maintenance of a relatively stable environment inside an organism usually involving feedback regulation. Homeostasis is maintained in face of Deals with temperature, pH, chemical concentrations,pressure, oxygen levels. Occurs through negative feedback loops. Various forms: simple thermostat in house turns off heater when above a certain temperature and on when below a certain temperature Involves stimulus, sensor,integrating center, effector and response. More efficient control has two sensors and two effectors. Can be antagonistic to each other, such as, one cools, the other heats. Precise control through proportional control, not all-or-none, furnace comes on a little bit if the house a bit cold. Examples in humans: vasoconstriction, change in metabolic rate, shivering. Physiological responses for high body temp: blood goes to body surface, sweating, behavioral changes (get out of sun). Positive feedback loop: effector increases deviation from set point. Amplifies reaction. Like blood clotting process, uterine contraction during childbirth. Negative feedback must exist at some point for control. Osmotic environments and regulations. Hypoosmotic: having less osmotic potential than nearby fluid Hyperosmotic: having more osmotic potential than nearby fluid Isoosmotic: having equal osmotic potential than nearby fluid Glomerulus: reduces volume of kidney Fish started in salt water, spread to fresh water, later reinvaded salt-water environment. Terrestrial animal water sources:&lt;br&gt; 1. drinking&lt;br&gt; 2. moist foods&lt;br&gt; 3. from breakdown of metabolic molecules like fats. (Desert kangaroo rats get 90% of their water from metabolism.) Secretion of nitrogenous wastes: from metabolism of amino acids, amino group has to be removed in one of three basically interchangeable chemical forms:&lt;br&gt; 1. ammonia (aquatic life)&lt;br&gt; 2. urea (mammals)&lt;br&gt; 3. uric acid (birds) Ammonia very toxic, soluble, and cheap to produce. Easy to expel for bony fishes. Urea: low toxicity, good solubility, more costly to lose as it contains other groups on it. Must be released in solution, water cost. Uric acid (white part of bird poo) low toxicity, insoluble, secreted with little water loss, more costly side groups lost than the others. Mammalian kidney: Structure: fist-sized organ in lower back. About 1/5 of blood from aorta at any time is passing through kidneys. Blood passes through kidney many times a day. Nephron: structural and functional unit of kidney. Bowmans capsule: funnel-like opening, contains primary filter, the glomerulus. Proximal convoluted tubule: receives stuff from Bowmans capsule. Loop of Henle: descends and ascends. Vasa recta: capillaries that surround the Loop of Henle. Glomerulus: main filter of the nephron, located within the Bowman's capsule Kidney properties and processes important to its function 1. Active transport of solutes from one fluid to another against a concentration gradient, Na+ actively transported out of filtrate by cells of the thick ascending loop of Henley into the interstitial fluid&lt;br&gt; 2. Passive movement of solutes and water from one fluid to another(down a concentration gradient), movement of water and NaCl out of descending loop of Henley into interstitial fluid.&lt;br&gt; 3. Differential permeability of cells in different regions of the nephron to movement of water and solutes, ascending thick look is impermeable to water, descending portion is permeable to water&lt;br&gt; 4. Hormonal control of that permeability, antidiuretic hormone(ADH) increases permeability of collecting due to water, resulting in reduced volume of filtrate and thus more concentrated urine.&lt;br&gt; 5. Increasing solute concentration in the interstitial fluid of the kidney, from the cortex to the deepest medulla, maintained by a countercurrent multiplier mechanism 

=Epithelial tissue= Comes from various sources, ectodermal or endodermal material. Cell sheet lines a surface or body cavity. One side, called freesurface or Apical, is exposed to The other side rests on the basal layer. Epithelial tissue is not penetrated by blood vessels. Two categories: Classified on two features: Cell shape at free surface: Two specialized epithelia: 

=Connective tissue = This is a “grab bag” category of diverse tissue types. Functions include binding and supporting. Types include bone, cartilage, fibrous connective tissue, blood and adipose (fat) tissue. If you took away everything in the body except the connective tissue, you’d still be able to see the basic form of the body. Form: distinctive cells surrounded by a cell matrix made of extra-cellular fiber grounded in a ground substance (excluding blood) Types:&lt;br&gt; 1. connectile connective tissues (can be 1. loose or 2. dense)&lt;br&gt; 2. special connective tissue (includes blood, bones and cartilage). Fibroblasts form connective tissue proper;&lt;br&gt; chondoroblasts form cartilage;&lt;br&gt;asdasdasdasdasdasdads osteoblasts form bone;&lt;br&gt; and blood is formed from various sources. Ground substance: “unstructured” material that fills space between cells and contains fibers. Made of&lt;br&gt; 1. interstitial fluid (bathes cells)&lt;br&gt; 2. proteoglycans (protein core with attached polysaccharides, glycoaminoglycans or GAGs such as chondroitin sulfate, keratin sulfate, and hyalronic acid, whose consistency is syrupy to gelatin-like)&lt;br&gt; 3. cell-adhesion proteins (connect connective tissue cells to the fibers). Fibers of connective tissue:&lt;br&gt; 1. Collagen (flexible protein resistant to stretching, tensile strength, most abundant protein in animals, white)&lt;br&gt; 2. elastin (rubbery, resilient protein, in dermis, lungs, blood vessels, yellow when fresh)&lt;br&gt; 3. andreticulin (like collagen). Loose connective tissue: found beneath skin, anchors muscles,nerves etc. Include fibroblasts, macrophages, mast cells,and adipose cells. Fibers include collagen and elastic fibers. Ground substance is “syrupy”. Adipose included. Dense connective tissue: largely densely packed fibers of collagen or elastin regularly or irregularly arranged. Forms tendons and ligaments, coverings of muscles, capsules around organs and joints, and dermis of skin. Cartilage vs. bone Cartilage: There are three cartilage types:&lt;br&gt; 1. hyaline cartilage&lt;br&gt; 2. fibrocartilage (fibrous cartilage)&lt;br&gt; 3. elastic cartilage Hyaline cartilage: most widespread cartilage type, in adults forms articular surfaces of long bones, rib tips, rings of trachea, and parts of skull. Mostly collagen, name refers to glassy appearance. In embryo, bones form first as hyaline cartilage, later ossifies. Found in tracheal rings. Few collagen fibers. Fibrous cartilage: have lots of collagen fibers. Found in intervertebral discs, pubic symphesis. Grades into dense tendon and ligament tissue. Elastic cartilage: springy and elastic. Found in internal support of external ear and in epiglottis, yellow when fresh. Chondrocites (cartilage cells) rely on diffusion for nutrients, as cartilage has no direct blood supply, and no enervation (nerves). Can be loaded with calcium salts. Bone: Specialized connective tissue, calcium phosphate arranged in highly ordered unit called osteon, or Hyvercian system. Concentric rings around central canal with blood vessels and enervation (nerves). Bone varied, not all vertebrate bone is even cellular. Our concern: simple pattern for mammals. Lacuna (spaces in which osteocytes found); canaliculi (little canals) bigger diagonal cells, layers of bone called lamellae. Three types of bone cells, ending in&lt;br&gt; -blast, (mend bone)&lt;br&gt; -cyte (fortify bone) &lt;br&gt; -clast (tear down bone) Classified by&lt;br&gt; 1. appearance (spongy vs. hard)&lt;br&gt; 2. where found (outside or inside)&lt;br&gt; 3. how it is formed (endochondral cartilage model forms first and then is ossified, and entramembranous, bone forms directly without cartilage precursor) Example of endochondral bone formation: long bone begins to ossify from center shaft, calcified region expands and cuts off diffusion of nutrients as bone replaces cartilage. In young mammals, secondary ossification centers then form at bone ends, growth has stopped by sexual maturity as all primary bone is ossified. In other animals, bones continue growing throughout their lifetime. Three types of intramembrous bone:&lt;br&gt; 1. dermal bone&lt;br&gt; 2. sessamoid bone&lt;br&gt; 3. perichondral bone. Dermal bone forms skull, shoulder/pectoral girdle, and integument, descended from dermal armor of ancestor. Comes from mesoderm, in dermis of skin. Sessamoid bones: form directly in tendons. Example: kneecap, also in wrist. Deals with stress. Perichondral bone means “around cartilage,” forms around cartilage or bone. Functions in bone repair and in ossification of endochondral bone. Bone remodeling and repair: bone has mineral structure, and develops tiny fractures, which, under stress, can lead to larger fractures. To combat this, bone is constantly replaced. Osteoclasts channel through existing bone, tear down and leave behind osteoblasts and lacuna, leaving osteocytes. Continually resets mineral structure of bone, and is preventative maintenance. When bone broken, callus forms in open ends, periosteum gives rise to new bone with calcium and new bone matrix, leaves irregular mend. Later, osteoblasts continue fixing over time and slowly removing imperfection. 

=Muscle tissue = Mesodermal in origin, muscle has several functions: supply force for movement, restrain movement, proper posture, act on viscera (internal organs) for peristalsis (moving food down digestive tract), give body shape, form sphincters, (such as in esophagus, between stomach and intestine, large and small intestine, in anus), in sheets of muscles, affect air flow in and out of lungs, line blood vessels and play vital role in circulation. Secondary roles: heat production (shivering a specialized heat production to supplement metabolism). Muscles co-opted to other non-original functions: sharks detect electrical field created by fish muscles. Some fish formed electric organs, create current strong enough to repel predators or stun prey. Other fish can use field as “radar” to see things and communicate with other animals. (Evolved independently in different groups). Different classifications: by color, (red or white) location, nature of nervous system control (voluntary or involuntary), embryonic origin, or by general microscopic appearance (striated, smooth, and cardiac.) Striated muscle (or skeletal muscle): under voluntary control. Individual cells called fibers, grouped into fascicle. Myofibrils founding one cell made of even smaller myofilaments. Each striated cell very long and multi-nucleated. Fibers joined end to end to form longer composite fibers. Sarcomeres: repeating units make up myofibrils. Two kinds of myofilaments, thick kind made up of myosin and thin of actin. Striations visible in light microscope, smaller part only with electron microscope. Cardiac muscle: occurs only in heart. Light banding visible under light microscope. Each band short, principally mononucleate (occasionally dinucleate) often branched, joined together with intercollated discs. Involuntary. Waves of contraction spread through intercollated discs. Initiated by nerve stimulation or can originate in the heart itself (useful in heart transplants.) Smooth muscle: no striations visible with light microscope. Almost entirely visceral function: digestion, sphincters, urogenital tracts, piloerectory muscles (make hairs stand up), lungs. Non-voluntary control. Slow and sustained action. Each cell mononucleate, short, fusiform (spindly) in shape, cells usually uniform in size. Striated muscle contraction: Muscle broken into units called fascicles, in units of myofibrils. Repeating units called sarcomeres, consisting of two kinds of myofilaments: &lt;br&gt; 1. thick, myosin filament &lt;br&gt; 2. thin, actin filament. Sarcomere: Thick and thin filaments interspersed in ordered grid. Sliding filament theory: thick and thin filaments move past each other in opposite direction, shortening length. Longer muscles contract more rapidly than short ones (see cell bio for details). Myosin molecule: two polypeptides twisted together with two globular heads at end. Myosin filament: many slender myosin molecules together. Actin filament: chain of actin single, tropomyosin strands with repeated globular troponin, and with actin. All play role in muscle contraction. Myocin heads have sites that bind to actin. Actin filaments have many regular sites that can bind to myosin. Troponin has four sites: &lt;br&gt; 1. one to bind myosin &lt;br&gt; 2. one for actin &lt;br&gt; 3. one for tropomyocin &lt;br&gt; 4. one for calcium ions Nerve signal reaches muscle, triggers release of chemical signal called neurotransmitter, that diffuses across cell membrane (sarcolimic reticulum) and binds to receptors in it. Receptor is acetylcholine, ACH. When there is enough nerve signal, the message travels through t-line to sarcoplasmic reticulum to release calcium ions. Lacking calcium, tropomyosin site blocked. In calcium, myosin binding sites exposed and heads bind to actin molecules, delivering force to move fibers in relation to each other. Myocin head then interacts with ATP to get “recocked”, if myosin still exposed then it fires again and results in further muscle contration. If there is no further nerve signal, sarcoplasmic reticulum sequesters Ca+ ions again and no recocking occurs. Quirari (or curare): known from movies, used in South America, blocks acetylcholine receptors in cell and causes skeletal paralysis. Victim dies of asphyxiation because he can’t breathe. Duchenne's muscular dystrophy: degeneration of sarcolema, plasma membrane of muscle cell unable to release signal and quickly atrophies. Fast and slow twitch fibers: vertebrate muscle fiber. Terms relative within one group of animals. Differences related to differences in enervation, type of myocin, and actin activation. Two parts of force generated by muscle: 1. active component 2. elastic component (energy stored in muscle when stretched by gravity or another force. Stored in muscle elastic tissue around tendons. Especially important in limb oscillation, like running, or trunk twisting, like fish swimming. Up to 90% of stored elastic energy can be recovered.) How does a muscle match its power to its job? Two ways: &lt;br&gt; 1. rate modulation, derived from frequency of nervous stimulation of muscle, force increases as frequency of stimulation increases up to point of tetanus. &lt;br&gt; 2. selective involvement of motor units, a given neuron enervates a fixed number of muscle cells, (a motor unit), and force is increased by recruiting more motor units. Motor units may be small, such as in eye, or larger, like in leg muscle. How do muscles grow stronger? &lt;br&gt; 1. add more myofilaments, increases cross sectional area by up to 50%, more little ratchets working &lt;br&gt; 2. proliferation in blood vessels and connective tissue around muscle Muscle strength is relative to cross sectional area, not length. Not always feasible to add more cross sectional area. Pinnate fibers: oriented obliquely (Y-shaped) to minimize muscle mass, in certain circumstances, like calf muscle. Spreads muscle out. Velocity of shortening greater in long muscle than short. Why? Contraction tied to relation between fibers, and to total length of muscle. Both long and short muscles reach same percentage of contraction in same unit time, but distance covered by the longer muscle is greater. Synergist muscles: muscles work together to produce motion in same general direction. Bicep shares work with brachialis. Antagonist muscles: muscles that oppose each other. Bicep pulls forearm in, triceps pulls it back out. Origin vs. insertion: origin is the end of the muscle that more fixed in its attachment to the body. The more movable end called insertion. Fixators: muscles that act to stabilize a joint or lever system. Like upper arm when you clench your fist hard. Flexors and extensors: applied mainly to limbs. Flexor bends one part relative to another about limb, extensor straightens it. Adductor and abductor: adductor draws a limb toward the ventral surface. Abductor moves limb away from ventral surface. (Adduct: drawn toward; abduct: carry away). 

The word "stoichiometry" derives from two Greek words: "stoicheion" (meaning "element") and "metron" (meaning "measure"). Stoichiometry deals with calculations about the masses (sometimes volumes) of reactants and products involved in a chemical reaction. It is a very mathematical part of chemistry, so be prepared for lots of calculator use. Jeremias Benjaim Richter (1762-1807) was the first to lay down the principles of stoichiometry. In 1792 he wrote: "Die stöchyometrie (Stöchyometria) ist die Wissenschaft die quantitativen oder Massenverhältnisse zu messen, in welchen die chymischen Elemente gegen einander stehen." [Stoichiometry is the science of measuring the quantitative proportions or mass ratios in which chemical elements stand to one another.] Molar Calculations. Your Tool: Dimensional Analysis. Luckily, almost all of stoichiometry can be solved relatively easily using dimensional analysis. Dimensional analysis is just using units, instead of numbers or variables, to do math, usually to see how they cancel out. For instance, it is easy to see that: It is this principle that will guide you through solving most of the stoichiometry problems (chemical reaction problems) you will see in General Chemistry. Before you attempt to solve a problem, ask yourself: what do I have now? where am I going? As long as you know how many (units) per (other units), this will make stoichiometry significantly easier. Moles to Mass. "How heavy is 1.5 mol of lead? How many moles in 22.34g of water?" Calculating the mass of a sample from the number of moles it contains is quite simple. We use the molar mass (mass of one mole) of the substance to convert between mass and moles. When writing calculations, we denote the molar mass of a substance by an upper case "M" (e.g. M(Ne) means "the molar mass of neon"). As always, "n" stands for the number of moles and "m" indicates the mass of a substance. To find the solutions to the two questions we just asked, let's apply some dimensional analysis: Can you see how the units cancel to give you the answer you want? All you needed to know was that you had 1.5 mol Pb (lead), and that 1 mol Pb weighs 207.2 grams. Thus, multiplying 1.5 mol Pb by 207.2 g Pb and dividing by 1 mol Pb gives you 310.8 g Pb, your answer. Mass to Moles. But we had one more question: "How many moles in 22.34g of water?" This is just as easy: Where did the 18 g H2O come from? We looked at the periodic table and simply added up the atomic masses of two hydrogens and an oxygen to get the molecular weight of water. This turned out to be 18, and since all the masses on the periodic table are given with respect to 1 mole, we knew that 1 mol of water weighed 18 grams. This gave us the relationship above, which is really just (again) watching units cancel out! Calculating Molar Masses. Before we can do these types of calculations, we first have to know the molar mass. Fortunately, this is not difficult, as the molar mass is exactly the same as the atomic weight of an element. A table of atomic weights can be used to find the molar mass of elements (this information is often included in the periodic table). For example, the atomic weight of oxygen is 16.00 amu, so its molar mass is 16.00 g/mol. For species with more than one element, we simply add up the atomic weights of each element to obtain the molar mass of the compound. For example, sulfur trioxide gas is made up of sulfur and oxygen, whose atomic weights are 32.06 and 16.00 respectively. The procedure for more complex compounds is essentially the same. Aluminium carbonate, for example, contains aluminium, carbon, and oxygen. To find the molar mass, we have to be careful to find the "total" number of atoms of each element. Three carbonate ions each containing three oxygen atoms gives a total of nine oxygens. The atomic weights of aluminium and carbon are 26.98 and 12.01 respectively. Empirical Formulae. The empirical formula of a substance is "the simplest ratio of the number of moles of each element in a compound". The empirical formula is ambiguous, e.g. the formula CH could represent CH, C2H2, C3H3 etc. These latter formulae are called molecular formulae. It follows that the molecular formula is always a "whole number multiple" of the empirical formula for a compound. Calculating the empirical formula is easy if the relative amounts of each element in the compound are known. For example, if a sample contains 1.37 mol oxygen and 2.74 mol hydrogen, we can calculate the empirical formula. A good strategy to use is to divide all amounts given by the smallest non-integer amount, then multiply by whole numbers until the simplest ratio is found. We can make a table showing the successive ratios. The empirical formula of the compound is H2O. Here's another example. A sample of piperonal contains 1.384 mol carbon, 1.033 mol hydrogen and 0.519 mol oxygen. The empirical formula of piperonal is C8H6O3. Converting from Masses. Often, we are given the relative composition by mass of a substance and asked to find the empirical formula. These masses must first be converted to moles using the techniques outlined above. For example, a sample of ethanol contains 52.1% carbon, 13.2% hydrogen, and 34.7% oxygen by mass. Hypothetically, 100g of this substance will contain 52.1 g carbon, 13.2 g hydrogen and 34.7 g oxygen. Dividing these by their respective molar masses gives the amount in moles of each element (as we learned above). These are 4.34 mol, 13.1 mol, and 2.17 mol respectively. The empirical formula of ethanol is C2H6O. Molecular Formula. As mentioned above, the molecular formula for a substance equals the count of atoms of each type in a molecule. This is always a "whole number multiple" of the empirical formula. To calculate the molecular formula from the empirical formula, we need to know the molar mass of the substance. For example, the empirical formula for benzene is CH, and its molar mass is 78.12 g/mol. Divide the actual molar mass by the mass of the empirical formula, 13.02 g/mol, to determine the multiple of the empirical formula, "n". The molecular formula equals the empirical formula multiplied by "n". This shows that the molecular formula for benzene is 6 times the empirical formula of CH. The molecular formula for benzene is C6H6. Solving Mass-Mass Equations. A typical mass-mass equation will give you an amount in grams and ask for another answer in grams. For example, given the equation formula_1, find out how many grams of silver (Ag) will result from 43.0 grams of copper (Cu) reacting. Summary. To solve a stoichiometric problem, you need to know what you already have and what you want to find. Everything in between is basic algebra. In general, all you have to do is keep track of the units and how they cancel, and you will be on your way! 

An adverb is a word added to a verb, a participle, an adjective, or another adverb; and generally expresses time, place, degree, or manner: as, Adverbs can modify a verb, a clause, adjective or a phrase. Form. Adjectives are generally turned into adverbs with the addition of a "ly" suffix. "Ly" is a contraction of like, and is the most common termination of adverbs. When added to nouns, it forms adjectives; but a few of these are also used adverbially: as, "daily, weekly, monthly". Examples of adverbs are: In the first sentence, the adverb modifies the verb swimming. The adjective "quick" has had a "ly" added to it to make an adverb. In the second sentence, it modifies the entire sentence and in the final example, the adverb "much" modifies the adverb "faster". Comparative forms of adverbs. Adverbs have no modifications, except that a few are compared, after the manner of adjectives: as, "soon, sooner, soonest; long, longer, longest; fast, faster, fastest". The following are irregularly compared: "well, better, best; badly or ill, worse, worst; little, less, least; much, more, most; far, farther, farthest; forth, further, furthest". Kinds of adverbs. Adverbs may be reduced to four general classes; namely, adverbs of time, of place, of degree, and of manner. Besides these, it is proper to distinguish the particular class of conjunctive adverbs. Adverbs of time. Adverbs of time are those which answer to the question, when? how long? how soon? or how often? Of time present: as, "now, yet, today, nowadays, presently, instantly, immediately, straightway, directly, forthwith". Of time past: as, "already, just now, lately, recently, yesterday, formerly, anciently, once, heretofore, hitherto, since, till now, long ago, erewhile, erst". Of time to come: as, "tomorrow, hereafter, henceforth, henceforward, by-and-by, soon, erelong, shortly". Of time relative: as, "when, then, first, just, before, after, while, whilst, meanwhile, as, till, until, seasonably, betimes, early, late, whenever, afterward, afterwards, otherwhile, otherwhiles". Of time absolute: as, "always, ever, never, aye, eternally, forever, perpetually, continually, incessantly, endlessly, evermore, everlastingly". Of time repeated: as, "often, oft, again, occasionally, frequently, sometimes, seldom, rarely, daily, weekly, monthly, yearly, annually, once, twice, thrice, or three times". Above thrice, we use only the phrases four times, five times, six times, etc. Times, for repetitions, or instances, may be supposed a noun; but such phrases often appear to be used adverbially. Adverbs of degree. Adverbs of degree are those which answer to the question, how much? how little? or to the idea of more or less. Of excess or abundance: as, "much, more, most, too, very, greatly, far, besides; chiefly, principally, mainly, mostly, generally; entirely, full, fully, completely, perfectly, wholly, totally, altogether, all, quite, clear, stark; exceedingly, excessively, extravagantly, intolerably; immeasurably, inconceivably, infinitely". Of equality or sufficiency: as, "enough, sufficiently, competently, adequately, proportionally, equally, so, as, even, just, exactly, precisely". Of deficiency or abatement: as, "little, less, least, scarcely, hardly, scantly, scantily merely, barely, only, but, partly, partially, nearly, almost, well-nigh, not quite". Of quantity in the abstract: as, "how, however, howsoever, everso, something, anything, nothing, a groat, a sixpence", and other nouns of quantity used adverbially. Adverbs of manner. Adverbs of manner are those which answer to the question, how? or, by affirming, denying, or doubting, show how a subject is regarded. Of manner from quality: as, "well, ill, wisely, foolishly, justly, wickedly", and many others formed by adding "ly" to adjectives of quality. Of affirmation or assent: as, "yes, yea, ay, verily, truly, indeed, surely, certainly, doubtless, undoubtedly, assuredly, certes, forsooth, amen". Of negation: as, "no, nay, not, nowise, noway, noways, nohow". Of doubt or uncertainty: as, "perhaps, haply, possibly, perchance, peradventure, maybe". Of mode or way: as, "thus, so, how, somehow, nohow, anyhow, however, howsoever, like, else, otherwise, across, together, apart, asunder, namely, particularly, necessarily, hesitatingly, trippingly, extempore, headlong, lengthwise". Adverbs of place. Of place in which: as, "where, here, there, yonder, above, below, about, around, somewhere, anywhere, elsewhere, otherwhere, everywhere, nowhere, wherever, wheresoever, within, without, whereabout, whereabouts, hereabout, hereabouts, thereabout, thereabouts". Of place to which: as, "whither, hither, thither, in, up, down, back, forth, aside, ashore, abroad, aloft, home, homewards, inwards, upwards, downwards, backwards, forwards". Of place from which: as, "whence, hence, thence, away, out, off, far, remotely". Of the order of place: as, "first, secondly, thirdly, fourthly", etc. Thus, secondly means in the second place; "thirdly, in the third place"; etc. Conjunctive adverbs. The conjunctive adverbs are those which perform the office of conjunctions. The following words are the most frequently used as conjunctive adverbs: "after, again, also, as, before, besides, consequently, else, ere, even, furthermore, hence, how, however, moreover, nevertheless, as well, otherwise, since, so, still, till, then, thence, therefore, too, until, when, where, wherefore, whither, while". The adverbs of cause: "why, wherefore, therefore"; but the last two of these are often called conjunctions. The pronominal compounds: "herein, therein, wherein", etc. A short syntax. Adverbs relate to verbs, participles, adjectives, or other adverbs: as, "How blessed," except the following cases: independent adverbs, as "No," the word "amen", as "These things say the amen," an adverb before preposition, as "All along", and "much", "little", "far", and "all", as "Thus far is right." 

An adjective is a word added to a noun or pronoun, and generally expresses quality: as, A wise man; a new book; you two are diligent. Adjectives may be divided into six classes; namely, common, proper, numeral, pronominal, participial, and compound. A common adjective is any ordinary epithet, or adjective denoting quality or situation: as, "good, bad, peaceful, warlike, eastern, western, outer, inner". A proper adjective is an adjective formed from a proper name: as, "American, English, Platonic, Genoese". A numeral adjective is an adjective that expresses a definite number: as, "one, two, three, four, five, six", etc. A pronominal adjective is a definitive word which may either accompany its noun, or represent it understood: as, That is, A participial adjective is one that has the form of a participle, but differs from it by rejecting the idea of time: as, A compound adjective is one that consists of two or more words joined together, either by the hyphen or solidly: as, "nut-brown, laughter-loving, four-footed; threefold, lordlike, lovesick". Cardinal: "one, two, three, four, five, six, seven, eight, nine, ten, eleven, twelve, thirteen, fourteen, fifteen, sixteen, seventeen, eighteen, nineteen, twenty, twenty-one, twenty-two", etc. Ordinal: "first, second, third, fourth, fifth, sixth, seventh, eighth, ninth, tenth, eleventh, twelfth, thirteenth, fourteenth, fifteenth, sixteenth, seventeenth, eighteenth, nineteenth, twentieth, twenty-first, twenty-second", etc. Multiplicative: "single or alone, double or twofold, triple or threefold, quadruple or fourfold, quintuple or fivefold, sextuple or sixfold, septuple or sevenfold, octuple or eightfold", etc. All that occur above decuple or tenfold, are written with a hyphen, and are usually of round numbers only: as, "thirty-fold, sixty-fold, hundred-fold". Adjectives word order. In English, multiple adjectives should follow a certain pattern as follows: Although you generally shouldn't use more than 2 or 3 adjectives together it is possible to make a sentence like this: Modifications. Adjectives have, commonly, no modifications but the forms of comparison. Comparison is a variation of the adjective, to express quality in different degrees: as, "hard, harder, hardest; soft, softer, softest". There are three degrees of comparison; the positive, the comparative, and the superlative. The "'positive degree" is that which is expressed by the adjective in its simple form: as, The comparative degree is that which is more or less than something contrasted with it: as, The superlative degree is that which is most or least of all included with it: as, Those adjectives whose signification does not admit of different degrees, cannot be compared: as, "two, second, all, every, immortal, infinite". Those adjectives which may be varied in sense, but not in form, are compared by means of adverbs: as, "fruitful, more fruitful, most fruitful; fruitful, less fruitful, least fruitful". Regular comparison. Adjectives are regularly compared, when the comparative degree is expressed by adding "er", and the superlative, by adding "est" to them: as, "great, greater, greatest; mild, milder, mildest". In the variation of adjectives, final consonants are doubled, final "e" is omitted, and final "y" is changed to "i", agreeably to the rules for spelling: as, "hot, hotter, hottest; wide, wider, widest; happy, happier, happiest". The regular method of comparison belongs almost exclusively to monosyllables, with dissyllables ending in "w" or "y", and such others as receive it and still have but one syllable after the accent: as, "fierce, fiercer, fiercest; narrow, narrower, narrowest; gloomy, gloomier, gloomiest; serene, serener, serenest; noble, nobler, noblest; gentle, gentler, gentlest". Comparison by adverbs. The two degrees of superiority may also be expressed with precisely the same import as above, by prefixing to the adjective the adverbs "more" and "most": as, "wise, more wise, most wise; famous, more famous, most famous; amiable, more amiable, most amiable". The degrees of inferiority are expressed, in like manner, by the adverbs "less" and "least": as, "wise, less wise, least wise; famous, less famous, least famous; amiable, less amiable, least amiable". The regular method of comparison has, properly speaking, no degrees of this kind. Nearly all adjectives that admit of different degrees, may be compared by means of the adverbs; but, for short words, the regular method is generally preferable: as, "quick, quicker, quickest"; rather than, "quick, more quick, most quick". Irregular comparison. The following adjectives are compared irregularly: "good, better, best; bad, evil, or ill, worse, worst; little, less, least; much, more, most; many, more, most". A short syntax. Adjectives relate to nouns or pronouns, as "Worldly enjoyments," except the following cases: an intervening verb, as "To err is human," arithmetical numbers, as "Four hundred and fifty-six men," an abstract adjective, as "Being sublime," and an adjective as abstract noun, as "Sensations of sublime." An adjective is placed immediately before noun, as "Vain man," except the following cases: pronouns, as "They left me weary," other words, as "A mind conscious of right," an action, as "Virtue renders life happy," admiration, as "Goodness infinite," a verb, as "Truth stands independent," a prefix "a", as afraid, the nature of a participle, as "The time then present," poetry, as "Isles atlantic," technical usage, as "Notary public," an adjective, as "A being infinitely wise," several adjectives, as "A woman, modest, sensible, and virtuous," empathy, as "Weighty is the anger," an adjective in predicate, as "We call the boy good," and an adjective as adverb, as "Particularly". 

A pronoun is a word used instead of a noun: as, Pronouns are not a requirement of a sentence, and it is possible for them to never to be used in sentences. However, many sentences become unwieldy without them: The pronouns in English language are twenty-four; and their variations are thirty-two: so that the number of words of this class, is fifty-six. Pronouns are divided into three classes; personal, relative, and interrogative. Pronouns also change depending on whether they refer to one person or thing (singular) or a group of people or things (plural). Personal pronouns. A personal pronoun or personal is a pronoun that shows, by its form, of what person it is: as, The simple personal pronouns are five: namely, "I", of the first person; "you (thou)", of the second person; "he, she", and "it", of the third person. The compound personal pronouns are also five: namely, "myself", of the first person; "yourself (thyself)", of the second person; "himself, herself", and "itself", of the third person. First person pronouns are used when referring to oneself: Second person pronouns are used to refer to someone who you are conversing with, the person the sentence is intended to be heard by: Third person pronouns are used when referring to something else that is outside the conversation, either some other person, or an object not capable of understanding or communicating: Third person singular pronouns are the only pronouns marked for gender. If gender is unknown, use "he or she" or use a plural. Relative pronouns. A relative pronoun or relative is a pronoun that represents an antecedent word or phrase, and connects different clauses of a sentence: as, The relative pronouns are "who, which, what, that, as", and the compounds "whoever or whosoever, whichever or whichsoever, whatever or whatsoever". "What" is a kind of double relative, equivalent to "that which" or "those which"; and is to be parsed first as antecedent, and then as relative: as, Interrogative pronouns. An interrogative pronoun or interrogative is a pronoun with which a question is asked: as, The interrogative pronouns are "who, which", and "what"; being the same in form as relatives. "Who" demands a person's name; "which", that a person or thing be distinguished from others; "what", the name of a thing, or a person's occupation and character. Pronouns have the same modifications as nouns; namely, persons, numbers, genders, and cases. Definitions universally applicable have already been given of all these things; it is therefore unnecessary to define them again in this place. The declension of a pronoun is a regular arrangement of its numbers and cases. Simple personals. The simple personal pronouns are thus declined: Compound personals. The word "self", added to the simple personal pronouns, forms the class of compound personal pronouns; which are used when an action reverts upon the agent, and also when some persons are to be distinguished from others. They all want the possessive case, and are alike in the nominative and objective. Thus: Relatives and interrogatives. The relative and the interrogative pronouns are thus declined: Compound relatives. The compound relative pronouns, "whoever or whosoever, whichever or whichsoever", and "whatever or whatsoever" are declined in the same manner as the simples. Thus: Unclear Usage of Pronouns. Although helpful to eliminate repetitiveness of nouns, pronouns, when used too much, can make a sentence extremely vague: as, The reader does not know what "it" is. The reader does not know who "they" are. Y'all. The pronoun "y'all" is a contraction of "You all". It is traditionally used in the south of the United States, where in the north you all is more common. "Y'all" follows the same conjugation rules as "they". Very often it is incorrectly spelled "ya'll". A short syntax. A pronoun must agree with its antecedent, as "This is the book; it is excellent," except the following cases: something indefinite, as "Tell me who it was," a neuter pronoun, as "I cannot view it," the pronoun "it", as "It is not for kings," the adjective "many", as "Many a great genius, they have no friends," enallage, as "We shall close our remarks," another sense, as "Lamps is of the plural number," nominatives, as "Who are you?", absolute nominatives, as "It need not be any wonder," possessives, as "Him whose yoke is easy," objectives, as "Those whom she persuaded," neuter verbs, as "Whom did you suppose me to be?", familiar language, as "The man [whom] I trust," omission of the relative, as "The worst thing [that] could happen," a collective noun, as "The council were divided," the conjunction "or", as "James or John will favour us with his company," the conjunction "and", as "Saul and Jonathan were pleasant in their lives," one person or thing, as "This great philosopher and statesman," empathy, as "The good man, and the sinner too, shall have his reward," and "each", "every", or "no", as "Every plant and every tree produces others after its kind." 

Chapter 14 Chapter 14. Pinophyta. The Division Pinophyta in the Kingdom Plantae comprises those species of plants that were formerly classified as the "modern" gymnosperms of the Class Coniferales—that is the "conifers". Unlike many of the gymnosperm groups covered in the previous chapter, "pinophytes" are today still broadly represented on the landscape, forming extensive forests in both the Northern and Southern hemispheres. This does not mean this is a recent group in the paleontological record. Pinophytes are found as fossils as far back as the Upper Carboniferous (Paleozoic Era). Thus, it is a very ancient group, but one still having significant ecological importance on the planet. Pinophytes are mostly evergreen trees (some are shrubs) and many have great commercial value for their wood. These terms, "gymnosperm" and "conifer" no longer have standing in modern taxonomic treatments of plants. However, both terms are still widely used, so you should have a grasp of what they mean, and how they fit into the taxonomic terms that have replaced them. Evolution of the Gymnosperms. Gymnosperms are very different from the earliest vascular plants. Gymnosperms have very reduced gametophyte generations- female ovules and male pollen. Many gymnosperms, like the Pinophytes, have secondary growth from a vascular cambium. While tree ferns are tall and have a trunk resembling a tree, they are only very superficially similar. They do not have woody growth the way trees with secondary growth from a vascular cambium do. A single mutation of a fern or a horsetail could not produce a functioning seed plant. So gymnosperms must be descended from some progymnosperm ancestor that evolved these adaptations from the ferns, but was not fit enough to remain on earth. One possible example of a progymnosperm is the spore bearing woody tree , which was at one point probably prevalent on earth. Laboratory Exercises for Chapter 14 »&lt;br&gt; 

unit 1 The following table shows the Spanish words for the four basic arithmetics operations: For the equals sign (=) you use es if the answer is 1, otherwise you use son: Problemas de matemáticas - "Math Problems". Write out as words and then solve the following math problems: 1.3*3 2. "10 / 2": ______________________________________ 3. "23 + 42": ________________________________________ 4. "99 - 9": ________________________________________ 5. "13 * 7": ________________________________________ 6. "75 / 3": ______ __________________________________ 7. "17 + 4": ________________________________________ 8. "67 - 66": ________________________________________ 9. "4 * 25": ________________________________________ 10. "78 / 13": _______________________________________ Soluciones a los ejercicios "Solutions to exercices" ^Lesson 2^ 

^ Indonesian ^ | « Lesson 7: Introducing Yourself | Lesson 8: My Family |Lesson 9: My Home » Bacaan "(Reading Comprehension)". Nama saya Mira. Saya tinggal di Bae bersama ayah dan ibu. Saya punya seorang kakak laki-laki. Namanya Anto. Saya juga punya adik perempuan. Namanya Wati. Kakak saya dan adik perempuan saya juga tinggal bersama kami. Ayah kerja di kantor setiap hari. Ibu tinggal di rumah, memasak makanan untuk kami. Setiap hari, saya, kakak dan adik perempuan saya pergi ke sekolah Terjemahan "(Translation):" My name is Mira I live in Bae with my father and mother. I have an older brother. His name is Anto. I also have a younger sister. Her name is Wati. My older brother and younger sister also live with us. Father works at the office everyday. Mother stays at home, cooking food for us. Every day, myself, my older and younger siblings go to school. Kosa Kata "(Vocabulary)". Note: The word "tinggal" is a tricky word. If it is used alone (not conjugated), it means to live or to stay. However, if you conjugate it, it means to leave or to die. Many people got confused with this. Think of it from an Indonesian cultural perspective. Many indigenous Indonesians believe that when a person dies, the soul is here to stay. Grammar. Awalan Me- "(Me- prefix)". In the passage, you'll notice the first prefix in Indonesian: "me-". It's the most important and commonly used in Indonesian. When it's combined with verbs like above (masak → memasak), it means the same as the infinitive form. The only thing is that we emphasize that now the verb is in active form. Almost all verb can be conjugated using "me-", but not all. Unfortunately, in order to know which verbs can go with "me-" , you must read a lot. The general rule of thumb is that if the verb is reflexive (i.e. doing it to ourself), it usually can't be conjugated with "me-". Even more so, don't think that the sense of reflexivity is the same to that of your language. Below is some words that cannot be conjugated with "me-", unless the meaning changes completely differently: The prefix "me-" can also apply to other type of words, such as noun and adjectives. However, the rough goal is still the same: To make an active verb. So, you can "verbize" a noun or an adjective. Details on these will be covered later. Note also that when words are conjugated with "me-", the spelling is changed a little bit. The spelling change is called "inflection". The inflection depends solely on the first letter of the original word. The rule on how the spelling changes can be viewed here: Prefix me- chapter. Akhiran -an "(Suffix -an)". Suffix "-an" is to "nominalize" a verb. Note the example from the passage: makan → makanan. Analoguous to the prefix "me-", you can virtually "nominalize" almost any verb you want using "-an" suffix, as long as it makes sense. Note, however, that the nominalisation of the verb does not correspond with the gerundive form (-ing form) expressing the "action" or "state of action" like in case of Western languages (i.e. eating), but the "target" or "object" of said action (i.e. food). You can think of it as "things that do what the verb says". Some of the verbs may also function as nouns already. For example: "tidur" (= sleep), which can be noun and verb at the same time. In this case, you cannot add "-an" suffix on it to make a noun out of it. (Note that the word "tiduran" does exist, but the meaning is "to lie down casually", not "sleep"). Note also that some words may have dual meanings, like "bangun", which may mean "to wake up" or "to build". Of course, if you add "-an" suffix, the second meaning is taken, i.e.: TODO: Load the chapter on suffix -an. Awalan Se- "(Prefix se-)". Prefix "se-" combined with noun would mean "one of that noun". In the example above, "orang" means "person". Therefore, "seorang" means "one person". In English, "se + noun" usually translates into "a / an". The "se + noun" word compounds are often used as a measure word, akin to those in Chinese. The measure words can be pretty complex. However, as being mentioned in the previous chapters, "sebuah" should be fine for most of the things. You must use "seorang" to indicate that the noun is a person, and use "seekor" for animals. TODO: Load the chapter on prefix se-. Grammar Summary. A simplified grammar summary for this chapter: &lt;hr&gt; ^ Indonesian ^ | « Lesson 7: Introducing Yourself | Lesson 8: My Family |Lesson 9: My Home » 

Lektion Vier Lesestück 5-1 ~ Eine Geschichte über St. Pölten. &lt;br&gt; Niederösterreich ist sowohl flächenmäßig als auch nach Einwohnern das größte der neun österreichischen Bundesländer. Sankt Pölten ist die Landeshauptstadt von Niederösterreich. Der Name St. Pölten geht auf den heiligen Hippolytos zurück, nach dem die Stadt benannt wurde. Die Altstadt befindet sich dort, wo vom 2. bis zum 4. Jahrhundert die Römerstadt "Aelium Cetium" stand. 799 wurde der Ort als "Treisma" erwähnt. Das Marktrecht erhielt St. Pölten um 1050, zur Stadt erhoben wurde es 1159. Bis 1494 stand St. Pölten im Besitz des Bistums Passau, dann wurde es landesfürstliches Eigentum. Bereits 771 findet sich ein Benediktinerkloster, ab 1081 gab es Augustiner-Chorherren, 1784 wurde deren Kollegiatsstift aufgehoben, das Gebäude dient seit 1785 als Bischofssitz. Zur Landeshauptstadt von Niederösterreich wurde St. Pölten mit Landtagsbeschluss vom 10. Juli 1986, seit 1997 ist es Sitz der Niederösterreichischen Landesregierung. &lt;br&gt; Vokabeln 5B.  Die Altstadt old town  Der Augustiner Augustinian  Der Besitz possession, holding  Das Bistum diocese  Der Bischofssitz bishop's see (a seat of a bishop's authority)  Die Bundesländer federal states  Die Chorherren men's choir  Das Eigentum proprietorship  Die Einwohner inhabitants  Das Gebäude premises  Die Geschichte history  Das Jahrhundert century  Das Kloster monastery, friary  Das Kollegiatsstift monastery college  Die Landeshauptstadt regional or state capital city  Die Landesregierung provincial (state) government  Der Landtagsbeschluss day of jurisdictional reorganization  Das Marktrecht right to hold markets  Der Name name  Der Ort place, spot, city  Die Römerstadt Roman town  Der Sitz official place  Bistum Passau a dioecian region in Bavaria  sowohl... als auch both... and  zurück auf goes back to  aufheben (hob auf, aufgehoben) merged in (or turned into?)  befinden sich situated, located  (befand sich, haben sich befunden)  finden sich* found (located)  benennen (benannte, benannt) call (as to label)  erhalten (erhielt, erhalten) receive  erheben (erhob, erhoben) arise, raise  erwähnen (erwähnte, erwähnt) mention  stehen (stand, gestanden) stand (stood, stood)  werden (wurde, [ist]geworden) become  ab from  auf up  bereits already  bis until, by, up to  flächenmäßig (no direct translation) ~ when measured in surface  heilig holy  landesfürstlich baronial or princely (holdings)  nach in terms of  um around 

Classical Latin is pronounced in a highly phonetic manner. Diphthongs: Digraphs: b, d, f, h, k, l, m, n, p, q, s, t, x, and z are pronounced as they are normally pronounced in English. 

The Domain Name System, most often known as simply DNS, is a core feature of the Internet. It is a distributed database that handles the mapping between host names (domain names), which are more convenient for humans, and the numerical . For example, "www.wikipedia.org" is a domain name and "130.94.122.199" the corresponding numerical internet address. The domain name system acts much like an automated phone book, so you can "call" "www.wikipedia.org" instead of "130.94.122.199". So, it converts human-friendly names such as "www.wikipedia.org" into computer-friendly (IP) addresses such as 130.94.122.199. It can also handle the reverse mapping, meaning that we can query for a name for "130.94.122.199", that return "larousse.wikipedia.org" DNS was first invented in 1983 by ; the original specifications are described in 882. In 1987 RFC 1034 and RFC 1035 were published which updated the DNS specification and made RFC 882 and RFC 883 obsolete. Subsequent to that there have been quite a few RFCs published that propose various extensions to the core protocols. DNS implements a hierarchical name space by allowing name service for parts of a name space known as zones to be "delegated" by a to subsidiary name-servers. DNS also provides additional information, such as alias names for systems, contact information, and which hosts act as mail hubs for groups of systems or domains. The present restriction on the length of domain names is 63 characters, excluding the "www." and ".com" or other extension. Domain names are also limited to a subset of characters, preventing many languages from representing their names and words correctly. The -based system, which maps strings into the valid DNS character set, has been approved by and adopted by some as a workaround. The DNS system is run by various flavors of DNS software, including: Any computer network can use DNS to implement its own private name system. However, the term "domain name" is most commonly used to refer to domain names implemented in the public DNS system. This is based on thirteen "root servers" worldwide, all but three of which are in the . From these thirteen , the rest of the Internet DNS name space is delegated to other DNS servers which serve names within specific parts of the DNS name space. An 'owner' of a domain name can be found by looking in the database: for most TLDs a basic WHOIS is held by ICANN, with the detailed WHOIS maintained by the domain registry which controlls that domain. For the 240+ Country Code TLDs the position is usually that the registry holds the entire authorative WHOIS for that extension, as part of their many functions. The current way the main DNS system is controlled is often criticized. The most common problems pointed at are that it is abused by monopolies or near-monopolies such as Inc., and problems with assignment of s. Some also allege that many implementations of DNS server software fail to work gracefully with dynamically allocated IP addresses, although that is the failure of specific implementations and not failures of the protocol itself. DNS uses and 53. Most DNS queries (such as name resolution requests) use UDP connections as the amount of data transferred is small and the session establishment overhead would introduce unnecessary traffic and load on nameservers. DNS zone file transfers between nameserver peers use TCP connections as the volume of data transferred is potentially much larger. A DNS domain definition (sometimes referred to as a 'zone file') consists of individual DNS records. There are several record types in common usage: Virtually all modern operating systems and network applications contain resolved libraries or routines for interrogating DNS services. However, OSs generally provide a command line interface for querying DNS servers. The Windows NT family of operating systems provides the 'nslookup' command. Unix-based operating systems may also offer 'nslookup' or 'dig' tools. nslookup can either be used interactively, or non-interactively. An example of non-interactive usage follows. In this example, we gather the A record for www.wikipedia.org from the client's default nameserver: "nslookup www.wikipedia.org" Nslookup is somewhat more powerful when used interactively. An example of this follows. In the example, we find the mail servers for the domain wikipedia.org: "nslookup" "&gt; set q=MX" "&gt; wikipedia.org" "Non-authoritative answer:" "wikipedia.org MX preference = 50, mail exchanger = mormo.org" "wikipedia.org MX preference = 10, mail exchanger = mail.wikimedia.org" See also: , , , 

Q1 The white areas are left neutral. Q2 When the electrons absorb the light, they get a little bit of extra energy. It is this increase in energy that allows them to escape to earth. Q3 The toner is given a negative charge so that it is attracted to the positively charged plate. Q4 The paper attracts the toner because it is positively charged. Note that the charge on the paper is stronger than the charge on the plate, so although the toner is still attracted to the plate, it is attracted to the paper even more. Q5 negative Q6 less paint falls on the floor becaue more is being attracted to the car. Remember: Unlike charges attract. «Back to Uses of Static Electricity 

What does your typical day looks like? Actions. Hindi uses two forms of address. One is honoring and the other is common. Honoring address is used when talking to first acquaintances, unfamiliar people, and to elders (relations and age)... Examples include The common type of address is used with friends and those of the same age group with whom one is familiar... Examples include For Hindi learners, it is better to stick to the first type of address. Hindi speakers normally use 'namaste' or 'namaskar' for all occasions (meaning good morning/afternoon/day/evening/night) especially with honorific use. Although these may also be used in pure Hindi: Phrases. नमस्ते (namaste) - a standard greeting (literally means "I bow to you") नमस्कार (namaskār) - a standard greeting आप कैसे हैं? (Āp kaise hain?) - How are you? आप का नाम क्या है? (Āp kā nām kyā hai?) - What is your name? आप कहाँ जा रहे हैं? (Āp kahān jā rahe hain?) - Where are you going? शुक्रिया (shukriyā) - thanks बहुत बहुत शुक्रिया (bahut bahut shukriyā) - thanks a lot आपका बहुत बहुत शुक्रिया (āpka bahut bahut shukriyā) - thank you so much प्रणाम (pranām) - a greeting, similar to नमस्ते (namaste) but used for much older or respectful people सही (sahī) - nice, right, good बहुत सही (bahut sahī) - very good आप की उम्र (āp kī umr) - your age आप का नाम (āp kā nām) - your name आप का पता (āp kā patā) - your address आज का मौसम कैसा है? (Āj ka mausam kaisā hai?) - How is the weather today? बैठिये (bathiye) - have a seat और सबज़ी लेंगे? (Aur sabzī lenge?) - Will you have more vegetables? और चावल लीजिये। (Aur chāwal lījiye.) - Have more rice. थोड़ी देर इन्तज़ार कीजिये। (Thoṛī der intzār kījiye.) - Wait for some time. ठहरिये (ṭhahriye) - wait a little समय क्या है? (Samay kyā hai?) - What's the time? Kya hal he apka Conversation. Here's a complete conversation in Hindi: नीरजः मैं विद्यालय /पाठशाला /स्कूल जा रहा हूँ। (Nīraj: Mainn vidhyalay/paathshaala/skūl jā rahā hūn.) "Neeraj: I am going to school." माँः ठीक है बेटा। (Mān: Ṭīk hai beṭā.) "Mom: All right son." नीरजः आज डिब्बे में क्या दिया है? (Nīraj: Āj dibbe men kyā diyā hai?) "What have you given me in tiffin?" माँः आज मैं ने सैंडविच और कुछ लड्डू दिये हैं। (Mān: Āj main ne sainḍwic aur kuch ladū diye hain.) "I have given some sandwich and sweet called ladoo." नीरजः वाह! (Nīraj: Wāh!) "Nice!" माँः बेटा समय से वापस आ जाना। (Mān: Beṭā samay se wāpas ā jānā.) "Son, come back on time." or "Honey, come back promptly." नीरजः हाँ माता जी! स्कूल ख़त्म होते सीधा घर आऊँगा। (Nīraj: Hān Mātā-jī! Skūl ḫatm hote sidhā ghar āūngā.) "Yes Mom! I'll be back home straight after school ends." 

Norsk ~ English Learning the Norwegian Language 

=Newton's Second Law= Momentum. A motion of a body is described by its momentum. Momentum is the product of a body's mass and its velocity. Momentum has both "magnitude" and "direction". Newton showed that bodies only change their motion when a force is applied, Newtons first law. The change in the momentum is equal to the applied force. Derivation of Newton's second law. Consider a snooker ball of mass "m" moving with a velocity "v"1. The ball hits the cushion and bounces off, it now has a velocity "v"2. The change in momentum is "mv"1 − "mv"2 = "m"("v"1 − "v"2). The change in velocity is the acceleration, and so the change in momentum and hence force "F" is mass times acceleration, "F" = "ma". Thus Action and Reaction are equal but opposite in direction. 

GCSE Science/Forces Pushing a car. Have you ever helped to push start a car? The hardest part is getting the car to move, but once it is rolling it takes much less effort, (especially if it is on a hill rolling down). To change the velocity of an object we have to do work on it.This means we have to supply energy to it. Once an object is moving it requires no further work, to keep it moving. Objects only slow down in everyday life because of friction forces which do work on the object. To do work on an object you have to apply a force to it, the amount of work you do depends on the acceleration of the object. Consider pushing a car quickly versus slowly. Pushing a car quickly requires more work. The rate of change in an objects velocity with respect to time is its acceleration. When you do work to push start a car, your energy is transferred to the kinetic energy of the car. When objects go faster they have more kinetic energy. We can slow an object down by removing energy from it. For instance if we use a spinning wheel to drive machinery, such as a generator, it will slow down as the energy is transferred to the generator and converted to electricity. 

The majority of the modules making up this book are based on notes very generously donated by Paul Doerder, Ph.D. and Ralph Gibson, Ph.D. both currently of the Cleveland State University. The book was initiated by , who donated many of his own class notes for other modules, and who is fleshing out the outline format of Dr. Doerder's notes into text. Users.  MD/PhD recent grad, enjoys teaching, hopes to help work up the Nervous System Tissue section. 

Notes on how to use slang. Foreign speakers. It is important to note that, as a foreigner, your use of slang will often be received as cute or funny, depending greatly upon your overall fluency in spoken French. To understand this, think about how it would sound to you if a foreigner—with a strong accent and odd rhythm of speech—came up to you and said "Dude, what a sketchy-ass hater that bizz-natch was, I totally was just like 'fuck off fo-sheezy'". Therefore, no matter how much slang you use in your native language, limiting your use of slang in French (proportionally to your level of fluency) will also limit how much you are patronised and giggled at by native listeners. Slang: consistency &amp; style. To use slang efficiently, it is important to maintain a consistency of style. Mixing styles might sound like saying: "Thy face, it is quite finely rawkin'". Translating 'fuck'. The English-language term 'fuck' is exceptional as it can serve as noun, verb, adjective, exclamation, and others. There is no such equivalent usage of any word in the French language, apart from 'putain', wich can be used as adjective and exclamation, "That's a fucking good car" = "C'est une putain de bonne voiture". Therefore, the translation of 'fuck' into French depends on the corresponding part of speech. Glossary.  Notes on Pronunciation:&lt;br&gt;*To feel how R should be pronounced, gargle with water, then try gargling without water.&lt;br&gt;That is what your throat should be doing when pronouncing the R.&lt;br&gt;*The U is hardest for English speakers. The back of the throat should be stretched out as if you see&lt;br&gt;a mouse and are saying "eee!", but the lips should be in a tight circle as if you are saying "ooo". exp. "à la con", stupid, in a stupid way. "J'ai cet examen à la con" = "I have this stupid test" Verlan. Verlan is roughly similar to English Pig Latin, in that certain words are split in half, and the two componenents switch positions, but do not necessarily retain all letters (due to French pronunciation patterns). For example, if you have word [12], in verlan it will become [2-1]. The word "verlan" is in itself an example of this; it comes from the word "l'envers" (meaning 'backwards'). Verlan is, unlike Pig Latin, quite commonly used among young adults and even adults. Common verlan expressions include: Common chat abbreviations. There are two general guidelines: 

Lektion Zwei für Fortgeschrittene Vokabeln 2-3.  das Delikatessengeschäft Deli, Delicatessen ("das Geschäft" = business)  der Hartkäse hard cheese  das Lebensmittel, die Lebensmittel food, foods  der Schmelzkäse soft cheese  die Schweinswurst pork sausage  der Schweizerkäse Emmenthaler cheese, Swiss cheese  das Stück piece  der Verkäufer sales clerk  das Würstchen small sausage  die Wurstsorten types of sausage  Bitte If you please  Nürnberger Schweinswürste a type of small, pork sausage (pl.)  finden find  heißen call, name  schmecken taste  suchen seek, look for  verkaufen sell (compare with "einkaufen" &amp; "der Verkäufer")  ähnlich similar  ein a, an, any, one  lecker tasty, delicious  nicht not  stückweise piecemeal, by the piece (compare with "das Stück") Grammatik 2-5 ~ Word Formation. As in any language, many words in German are constructed from other smaller words that provide similar meaning, although the connections can sometimes be obscured by the passage of time. Construction of new words from word combinations is especially prevalent with German nouns, and understanding word roots can therefore be helpful in learning new words. As an example, consider the phrase "Auf Wiedersehen" — the standard translation into English being 'Good bye', although it means literally 'upon reunion' (in essence, "until we meet again"). The noun, "das Wiedersehen", consists of "wieder", 'once again' (or 're-' as a prefix), and "sehen" or 'see'. The noun "die Geschäftsleute" provides a direct example of a compounded noun: the first part of each deriving from "das Geschäft" ('business') and the second part from "die Leute" ('people'). The gender of a compound noun follows the base or last noun. There are other examples in the this lesson, but these may not be immediately obvious unless you already have a good command of German words. However, you should train yourself to view new words in terms of the meanings of their component parts. Consider all of the various words used in this lesson to describe types of cheeses: "der Hartkäse", "der Schmelzkäse", "der Schweizerkäse"; or nouns and verbs related to buying and selling ("Kaufen und Verkaufen"). Grammatik 2-6 ~ Personal Pronouns: nominative case. Here are the personal pronouns in the nominative case: The nominative case is that of the subject of a verb. The pronoun subject of these sentences is underlined in the German and the English: This last sentence is an example from Gespräch 2-3 using the polite form of 'you'. Whether singular or plural must be established by context. This next sentence translates with "sie" as 'they': as evidenced by the form taken by the verb 'can' ("können"). Other uses of the nominative case in German will be explored in future lessons. Tables of the pronouns in all cases are summarized in the grammar appendix: Pronoun Tables. NOTE: An intransitive verb cannot be followed by an object in English or German. A pronoun following an intransitive verb such as 'to be' is called a predicate pronoun and should be in the nominative case. In English 'It is I' is correct; 'It is me' is incorrect. Grammatik 2-7 ~ More on verb forms. Just as English sometimes adds the verb "to be", forming the progressive, note also in Grammatik 2-2 (in both question sentence examples) that English also may insert the verb 'to do' (called the emphatic form), especially useful when forming a question. This is not done in German: Again, in the present tense, the English sentences: are all, in German: "Ich schreibe einen Brief." means: "Ich schreibe gerade einen Brief." Vokabeln 2-4.  der Brief letter  das Einkaufen shopping  der Finger, die Finger finger, fingers  das Kaufen buying (use of the verb form is preferred)  das Schwein pig (compare with "die Schweinswurst")  das Verkaufen selling  können can  schreiben write  jede any  zehn ten Andere Wörter 2A. Using these additional vocabulary words, you should be able to restate Gespräch 2-2 above, altering the meaning (or time of day) of the conversation.  der Abend evening  Guten Abend! Good Evening (greeting)  morgen früh tomorrow morning  zu wenig too little  abend evening  abends evenings  falsch false, wrong  morgen tomorrow  morgens in the morning  schlecht bad Übersetzung 2-2. Write these sentences in German. Pay attention to the additional words presented in "Andere Wörter 2A": 

Objective-C is an object-oriented programming language. It was named after the concept of adding objects to the C language. Objective C was introduced with NeXTSTEP and OPENSTEP, and was considerably extended in application due to its use with the Cocoa libraries under Mac OS X or the GNUstep libraries. However, you can program in Objective-C without these libraries if you wish. Because not everyone uses OPENSTEP or Mac OS X, we will introduce these library-specific details later. Familiarity with the is required, as Objective-C shares a lot with it, hence its name. See also.  __NOEDITSECTION__ 

The Normal Force. Why is it that we stay steady in our chairs when we sit down? According to the first law of motion, if an object is translationally in equilibrium (velocity is constant), the sum of all the forces acting on the object must be equal to zero. For a person sitting on a chair, it can thus be postulated that a normal force is present balancing the gravitational force that pulls the sitting person down. However, it should be noted that only some of the normal force can cancel the other forces to zero like in the case of a sitting person. In Physics, the term normal as a modifier of the force implies that this force is acting perpendicular to the surface at the point of contact of the two objects in question. Imagine a person leaning on a vertical wall. Since the person does not stumble or fall, he/she must be in equilibrium. Thus, the component of his/her weight along the horizontal is balanced or countered (opposite direction) by an equal amount of force -- this force is the "normal force" on the wall. So, on a slope, the normal force would not point upwards as on a horizontal surface but rather perpendicular to the slope surface. The normal force can be provided by any one of the four fundamental forces, but is typically provided by electromagnetism since microscopically, it is the repulsion of electrons that enables interaction between surfaces of matter. There is no easy way to calculate the normal force, other than by assuming first that there is a normal force acting on a body in contact with a surface (direction perpendicular to the surface). If the object is not accelerating (for the case of uniform circular motion, the object is accelerating) then somehow, the magnitude of the normal force can be solved. In most cases, the magnitude of the normal force can be solved together with other unknowns in a given problem. Sometimes, the problem does not warrant the knowledge of the normal force(s). It is in this regard that other formalisms (e.g. Lagrange method of undertermined coefficients) can be used to eventually "solve" the physical problem. Friction. When there is relative motion between two surfaces, there is a resistance to the motion. This force is called friction. Friction is the reason why people could not accept Newton's first law of Motion, that an object tends to keep its state of motion. Friction acts opposite to the direction of the original force. The frictional force is equal to the frictional coefficient times the normal force. Friction is caused due to attractive forces between the molecules near the surfaces of the objects. If two steel plates are made really flat and polished and cleaned and made to touch in a vacuum, it bonds together. It would look as if the steel was just one piece. The bonds are formed as in a normal steel piece. This is called cold welding. And this is the main cause of friction. The above equation is an empirical one--in general, the frictional coefficient is not constant. However, for a large variety of contact surfaces, there is a well characterized value. This kind of friction is called Coulomb friction. There is a separate coefficient for both static and kinetic friction. This is because once an object is pushed on, it will suddenly jerk once you apply enough force and it begins to move. Also, the frictional coefficient varies greatly depending on what two substances are in contact, and the temperature and smoothness of the two substances. For example, the frictional coefficients of glass on glass are very high. When you have similar materials, in most cases you don't have Coulomb friction. For static friction, the force of friction actually increases proportionally to the force applied, keeping the body immobile. Once, however, the force exceeds the maximum frictional force, the body will begin to move. The maximum frictional force is calculated as follows: The static frictional force is less than or equal to the coefficient of static friction times the normal force. Once the frictional force equals the coefficient of static friction times the normal force, the object will break away and begin to move. Once it is moving, the frictional force then obeys: The kinetic frictional force is equal to the coefficient of kinetic friction times the normal force. As stated before, this always opposes the direction of motion. Variables. Definition of Terms It's important to note, that in real life we often have to deal with viscose and turbulent friction - they appear when you move the body through the matter. Viscose friction is proportional to velocity and takes place at approximately low speeds. Turbulent friction is proportional to formula_1 and takes place at higher velocities. 

Before we begin As you gain programming experience you will appreciate the more specific explanations like The C Programming Wikibook. Introduction to programming. Programming has many uses. Some areas where programming and computer science in general are extremely important include artificial intelligence and statistics. Programming allows you to use computers flexibly and process data very quickly. When a program is written, it is written into a textual form that a human can understand. However, a computer doesn't directly understand what a human writes. It needs to be transformed into a way that the computer can directly understand. For example, a computer is like a person who reads and speaks German. You write and speak in English. The letter you write to the computer needs to be translated for the computer to speak. The program responsible for this work is referred to as the "compiler". You need to "compile" your English-like instructions, so that the computer can understand it. Once a program has been compiled, it is hard to "un-compile" it, or transform it back into English again. A programmer writes the program (to use our analogy, in English), called "source code", which is a human-readable definition of the program, and then the "compiler" translates this into "machine code". We recommend using the widely available gcc compiler. When we look at mathematical programming here, we will look at how we can write programs that will solve some difficult mathematical problems that would take us normally a lot of time to solve. For example, if we wanted to find an approximation to the root of the polynomial "x"5+"x"+1 - this is very difficult for a human to solve. However a computer can do this no sweat -- how? Programming language basics. We will be using the C programming language throughout the chapter, please learn about the basics of C by reading the first 7 lessons of "C Programming Tutorial" at About.com. Data Types. Size Constraints: Header files limits.h, and float.h Computers are machines based on Boolean logic. This means that the computer is based on some method of differentiating a state as true or false, or set and not set. Abstractly we think of computers as using 1's for true or set and 0's for false or not set. We refer to these 1's and 0's as bits. In computers we don't keep track of information as bits. Instead information in a computer is stored in addressable blocks called bytes. A byte is the smallest piece of memory that can be accessed in the computer that is not a bit. When we declare a [scalar] variable in a C program that memory has an address and a length. The address says where the memory starts, and the length states how many bytes are used to express the variable. The include file &lt;limits.h&gt; is used to define the size of addressable integer types and the include file &lt;float.h&gt; is used to define the size of addressable floating point types. The values in these files are compiler and computer dependent. This means that if you change compilers or compile your program on a different type of computer it may execute differently. Here are some of the values defined in these two files:  Exercises Programming With Integers. Discrete programming deals with integers and how they are manipulated using the computer. Understanding integer division. In C, the command will set aside some space in the computer memory, and we can refer to that space by the "variable" name number. In the computer's mind, number is an integer, nothing else. After numbers equals 1, not 1.5, this is due to that fact that / when applied to two integers will give only the integer part of the result. For example in C: If the number you are testing is between one and negative one - for example 2 / 5 or -2 / 5 then the result is undefined, although most compilers return 0. The "modular operator", %, returns the remainder resulting from integer division. For example in C: The sign of the result takes on the sign of the dividend as you would expect. For fractions that are between one and negative one the result is the same as the numerator.  Exercises Exercise 1 Write down your thoughts on what a program to explore division and modulus should do.&lt;br&gt; The following example will walk you through this exercise. Modeling Recursively defined functions. The factorial function n! is recursively defined: In C, if fact(n) is the functions as described above we want we should note that all recursively defined functions have a "terminating condition", it is the case where the function can give a direct answer to, e.g. fact(0) = 1. We can model the factorial functions easily with the following code and then execute it: The C function above models the factorial function very naturally. To test the results, we can compile the following code: We can also model the Fibonacci number function. Let fib(n) return the ("n" + 1)th Fibonacci number, the following should be clear we can model the above using C: Again, you shall see that modeling a recursive function is not hard, as it only involves translating the mathematics in C code. Modeling non-recursive functions. There are functions that involve only integers, and they can be modelled quite nicely using "functions" in C. The factorial function can be modeled quite simply using the following code  int n = 10; //get factorial for 10  int f = 1; //start f at 1  while(n &gt; 0) //keep looping if n bigger then 0  f = n * f; //f is now product of f and n  n = n - 1; //n is one less (repeat loop) Floating point Programming. Programs can not only be written with integer values, but also with various forms of floating-point values. You should normally use the "double" keyword to define a floating point number; the reason for this is that in many cases, the intuitive way to write an expression in floating point arithmetic is suboptimal. Floating point arithmetic is non-associative - in base-10, a system that has 2 places of accuracy has (1.0 + 0.02) + 0.04 = 1.0 (rounded down because 1.02 rounds to 1.0, and then 1.04 rounds down to 1.0), but 1.0 + (0.02 + 0.04) = 1.0 + 0.06 = 1.1. There also exists a "float" type, which uses 4 bytes instead of 8, but you should not use it unless you know what you are doing, since there is only 24 bits of accuracy, or roughly 9 base-10 significant digits. Beware: If you use 2 integer operands, it still performs integer arithmetic, so this prints 1, not 1.5 as you'd expect: An example of a definition of a floating-point number and then calculating 3/2: A caveat: floating point numbers do not perfectly represent all decimal numbers. Obviously, because the memory consumed by a variable is finite, it cannot represent an infinite number, but only an approximation of it. In addition, some numbers cannot be perfectly represented in floating point. In this code: the value of number is not 0.1, but actually 0.10000000000000001. An analogy of this limitation is the value of 1/3. In the decimal system, this value cannot be exactly represented with a finite number of 3s, as in 0.333... Since 2 and 5 are the only prime factors of 10 (the base of the decimal system), only fractions with denominators comprising products of 2s and 5s, such as 5/8 formula_4 or 231/250 formula_5 (but not 1/3 or 5/14) can be exactly represented by decimal numbers (0.625 and 0.924, respectively). Computers use base 2 arithmetic, so only fractions with denominators comprising products of 2s (powers of 2), such as 5/8 formula_4 or 231/256 formula_7 (but not 231/250, 1/3 or 1/10, as above) can be exactly represented by a floating point number. 

=Fields= A field is one of the more difficult concepts to grasp in physics. A field is an area or region in which an influence or force is effective regardless of the presence or absence of a material medium. Simply put, a field is a collection of vectors often representing the force an object "would" feel if it were placed at any particular point in space. With gravity, the field is measured in newtons, as it depends solely on the mass of an object, but with electricity, it is measured in newtons per coulomb, as the force on an electrical charge depends on the amount of that charge. Typically these fields are calculated based on canceling out the effect of a body in the point in space that the field is desired. As a result, a field is a vector, and as such, it can (and should) be added when calculating the field created by TWO objects at one point in space. Fields are typically illustrated through the use of what are called field lines or lines of force. Given a source that exerts a force on points around it, sample lines are drawn representing the direction of the field at points in space around the force-exerting source. There are three major categories of fields: Magnetism also has a field, measured in Tesla, and it also has field lines, but its use is more complicated than simple "force" fields. Secondly, it also only appears in a two-pole form, and as such, is difficult to calculate easily. The particles that form these magnetic fields and lines of force are called electrons and not magnetons. A magneton is a quantity in magnetism. 



Numbers. See ../Numbers from 0 to 999/ Colors. NOTE: "ah" in English represents the sound 'a' as in "far" 

Days. Days of the week. Note that as in most languages, but unlike English, days of the week and months are written with lower-case letters except when the word appears at the start of a sentence: "luni" not "Luni". 

Two-Dimensional Vectors. Introduction. In most mathematics courses up until this point, we deal with scalars. These are quantities which only need one number to express. For instance, the amount of gasoline used to drive to the grocery store is a scalar quantity because it only needs one number: 2 gallons. In this unit, we deal with vectors. A vector is a directed line segment -- that is, a line segment that points one direction or the other. As such, it has an initial point and a terminal point. The vector starts at the initial point and ends at the terminal point, and the vector points towards the terminal point. A vector is drawn as a line segment with an arrow at the terminal point: The same vector can be placed anywhere on the coordinate plane and still be the same vector -- the only two bits of information a vector represents are the magnitude and the direction. The magnitude is simply the length of the vector, and the direction is the angle at which it points. Since neither of these specify a starting or ending "location", the same vector can be placed anywhere. To illustrate, all of the line segments below can be defined as the vector with magnitude formula_1 and angle 45 degrees: It is customary, however, to place the vector with the initial point at the origin as indicated by the black vector. This is called the standard position. Component Form. In standard practice, we don't express vectors by listing the length and the direction. We instead use component form, which lists the height (rise) and width (run) of the vectors. It is written as follows: formula_2 Other ways of denoting a vector in component form include: formula_3 and formula_4 From the diagram we can now see the benefits of the standard position: the two numbers for the terminal point's coordinates are the same numbers for the vector's rise and run. Note that we named this vector formula_5 . Just as you can assign numbers to variables in algebra (usually formula_6), you can assign vectors to variables in calculus. The letters formula_7 are usually used, and either boldface or an arrow over the letter is used to identify it as a vector. When expressing a vector in component form, it is no longer obvious what the magnitude and direction are. Therefore, we have to perform some calculations to find the magnitude and direction. Magnitude. formula_8 where formula_9 is the width, or run, of the vector; formula_10 is the height, or rise, of the vector. You should recognize this formula as the Pythagorean theorem. It is -- the magnitude is the distance between the initial point and the terminal point. The magnitude of a vector can also be called the norm. Direction. formula_11 where formula_12 is the counter-clockwise angle made by the vector with the positive formula_13-axis. This formula is simply the tangent formula for right triangles. Vector Operations. For these definitions, assume: formula_14 Vector Addition. Vector Addition is often called "tip-to-tail" addition, because this makes it easier to remember. The sum of the vectors you are adding is called the resultant vector, and is the vector drawn from the initial point (tip) of the first vector to the terminal point (tail) of the second vector. Although they look like the arrows, the pointy bit is the tail, not the tip. (Imagine you were walking the direction the vector was pointing... you would start at the flat end (tip) and walk toward the pointy end.) It looks like this: Numerically: formula_15 Or more generally: formula_16 Scalar Multiplication. Graphically, multiplying a vector by a scalar changes only the magnitude of the vector by that same scalar. That is, multiplying a vector by 2 will "stretch" the vector to twice its original magnitude, keeping the direction the same. formula_17 Numerically, you calculate the resultant vector with this formula: formula_18 , where formula_19 is a constant scalar. As previously stated, the magnitude is changed by the same constant: formula_20 Since multiplying a vector by a constant results in a vector in the same direction, we can reason that two vectors are parallel if one is a constant multiple of the other -- that is, that formula_21 if formula_22 for some constant formula_19 . We can also divide by a non-zero scalar by instead multiplying by the reciprocal, as with dividing regular numbers: formula_24 Linear Functions. Given a function formula_25 that accepts a vector as input and returns a vector or scalar as the output, function formula_25 is considered to be "linear" if the following holds: More generally, when given a function formula_33 that has multiple vector valued parameters formula_34, function formula_33 is a "multi-linear" function if formula_33 is linear with respect to each parameter while holding all other parameters constant: For each formula_37 and vectors formula_38: If formula_44, then formula_33 is "bilinear". Bilinear functions include the dot product and the cross product. Dot Product. The dot product is a way of multiplying two vectors to produce a scalar value. Because it combines the components of two vectors to form a /scalar/, it is sometimes called a scalar product. If you were asked to take the 'dot product of two rectangular vectors' you would do the following: formula_46 It is very important to note that the dot product of two vectors does not result in another vector, it gives you a scalar, just a numerical value. Another common pitfall may arise if your vectors are not in rectangular ('cartesian') format. Sometimes, vectors are instead expressed in polar coordinates, where the first component is the vector's magnitude (length) and the second is the angle from the formula_13-axis at which the vector should be oriented. Dot products cannot be performed using the conventional method on these sorts of vectors; vectors in polar format must be converted to their equivalent rectangular form before you can work with them using the formula given above. A common way to convert to rectangular coordinates is to imagine that the vector was projected horizontally and vertically to form a right triangle. You could then use properties of sin and cos to find the length of the two legs the right triangle. The horizontal length would then be the x-component of the rectangular expression of the vector and the vertical length would be the y-component. Remember that if the vector is pointing down or to the left, the corresponding components would have to be negative to indicate that. With some rearrangement and trigonometric manipulation, we can see that the number that results from the dot product of two vectors is a surprising and useful identity: formula_48 where formula_12 is the angle between the two vectors. This provides a convenient way of finding the angle between two vectors: formula_50 Notice that the dot product is 'commutative', that is: formula_51 Also, the dot product of two vectors will be the length of the vector squared: formula_52 and by the Pythagorean theorem, formula_53 The dot product can be visualized as the length of a projection of one vector on to the other. In other words, the dot product asks 'how much magnitude of this vector is going in the direction of that vector?' Deriving the Dot Product. Start with the following definition for the dot product: formula_54 where formula_12 is the angle between formula_5 and formula_28. The formula formula_58 can be derived from the above definition through various approaches: Approach #1 One of the more direct approaches is to use the law of cosines. Create a triangle formula_59 with vertices formula_60, formula_61, and formula_62. The displacement formula_63, the displacement formula_64, and the displacement formula_65. The angle formula_66. The lengths of the sides of the triangle are formula_67, and formula_68, and formula_69. Applying the law of cosines gives: formula_70 formula_71 formula_72 formula_73 formula_74 Therefore formula_58. Approach #2 A more intuitive derivation of formula_58 uses the fact that the dot product is a bilinear operator. To establish that the dot product is a bilinear operator, the following must be established: Since it is readily apparent from the definition formula_54 that formula_84, linearity with respect to formula_28 implies linearity with respect to formula_5. It is hence only necessary to establish that the dot product is linear with respect to formula_28 to establish bilinearity. formula_88 where formula_89 is the "orthogonal projection" of vector formula_28 onto a line formula_91 whose direction is the direction of formula_5. formula_93 is the component of formula_28 that is parallel to formula_91. By drawing similar triangles, one observes that formula_93 is linear with respect to formula_28 while formula_5 is held constant. The dot product formula_99 is linear with respect to formula_28, and therefore the dot product is a bilinear operator. The bilinearity of the dot product now enables the derivation: formula_101 formula_102 formula_103 formula_104 formula_105 formula_106 Therefore formula_58. Applications of Scalar Multiplication and Dot Product. Unit Vectors. A unit vector is a vector with a magnitude of 1. The unit vector of u is a vector in the same direction as formula_5 , but with a magnitude of 1: The process of finding the unit vector of formula_5 is called normalization. As mentioned in scalar multiplication, multiplying a vector by constant formula_19 will result in the magnitude being multiplied by formula_19 . We know how to calculate the magnitude of formula_5 . We know that dividing a vector by a constant will divide the magnitude by that constant. Therefore, if that constant is the magnitude, dividing the vector by the magnitude will result in a unit vector in the same direction as formula_5 : formula_114 , where formula_115 is the unit vector of formula_5 Standard Unit Vectors. A special case of "Unit Vectors" are the "Standard Unit Vectors" formula_117 : formula_118 points one unit directly right in the formula_13 direction, and formula_120 points one unit directly up in the formula_121 direction: formula_122 formula_123 Using the scalar multiplication and vector addition rules, we can then express vectors in a different way: formula_124 If we work that equation out, it makes sense. Multiplying formula_13 by formula_118 will result in the vector formula_127 . Multiplying formula_121 by formula_120 will result in the vector formula_130 . Adding these two together will give us our original vector, formula_131 . Expressing vectors using formula_117 is called standard form. Projection and Decomposition of Vectors. Sometimes it is necessary to decompose a vector formula_5 into two components: one component parallel to a vector formula_134 , which we will call formula_135 ; and one component perpendicular to it, formula_136 . Since the length of formula_135 is formula_138 , it is straightforward to write down the formulas for formula_136 and formula_135 : formula_141 of which the top-most vector is the Normal vector, but the bottom half formula_142 is known as the curvature. Since the Normal vector points toward the inside of a curve, the sharper a turn, the Normal vector has a large magnitude, therefore the curvature has a small value, and is used as an index in civil engineering to reflect the sharpness of a curve (clover-leaf highways, for instance). The only other thing not mentioned is the Binormal that occurs in 3-d curves formula_143 , which is useful in creating planes parallel to the curve. 

Lisp is a programming language. It is named after the collapsed phrase List Processing. If you have programmed before and would like to see a little bit of how Lisp works and is different from other programming languages, you can get an overview. Dialects. Because Lisp itself is, technically, just seven operators, to become a useful language, much more needs to be implemented atop it. Common Lisp and Scheme are two such designs to create a useful programming language. Common Lisp is an ANSI standard, and features an extensive array of library functions. It is the more widely used of the two. Scheme is designed in a minimalistic fashion, with a very small amount of built in functions. This is probably true, but Scheme lacks many of the time-saving built-in functions of Common Lisp. Emacs Lisp is an implementation of Lisp in Emacs. newLISP is a Lisp-like, general-purpose scripting language See also. Lush - lisp-like object-oriented programming language 

See also. Finnish dialogues Finnish Language for Foreigners (English Translation) Finnish Language in Use (English Translation) Wikipedia articles about the Finnish language and the Finnish grammar. The Finnish language edition of Wikibooks. The Finnish language edition of Wikipedia. Discuss the Finnish language and all things Finnish. 

^^Contents^^ | «Introduction | Alphabet | Hello!» In short, the basic alphabet of the Finnish language: Finnish is a very phonetic language, as every pronunciation has its own letter. That is to say that things are "pronounced exactly as they are written" so SAMPA and IPA notations of Finnish words are almost identical to the written language. However, do not take this too literally; there are certainly many details in speech that cannot be easily expressed in written language, and Finnish is no exception. The glyphs 'ä' and 'ö' have been borrowed from Swedish. They are independent letters and phonemes (sounds), not modified nor accented letters. Changing 'Ä' into 'A' or 'Ö' into 'O' is akin to changing 'O' into 'Q'. Go through each of the letters of the alphabet and practice saying them: 



= Plural =  ^ Polish ^ Like most other , and unlike Chinese or Japanese, almost everything in Polish depends on number, which may be either singular or plural (pojedyncza and mnoga). Let's take a look at verbs first. And some nouns: As you can see, with the exception of person-masculine, the nominative and accusative are the same in the plural (in the person masculine accusatives are the same as genitives). Like in other Indo-European languages, the verb agrees with the subject.  ^ Polish ^ 

= Noun cases = There are seven cases of noun and adjective declension in Polish. So far we have only introduced the nominative ("mianownik") and the accusative ("biernik"), in . The whole list, in traditional order, is: Nominative. No need to memorize endings as nouns and modifiers are in this case. Use the basic, unmarked, dictionary forms. Accusative. Exception: The preposition "w" (in) takes the locative when indicating a state/condition, but to imply change you must use the preposition "do" + genitive. Instrumental. Exceptions: Most lone adjectives and "To jest X". Vocative. Only masculine and feminine singulars have a separate form. Before declension tables are provided, let's take a look at adjectives:  ^ Polish ^ 

= Adjectives = The adjective ("przymiotnik") is a very powerful part of speech in Polish. It declines very regularly depending on case, number and gender. It may be used as a noun. For example Polish adjectives "bogaty" (rich) and "czarny" (black) are very often used as "nouns" (taking forms "bogata" and "czarna" in feminine gender). While there is only one pattern, final consonant group takes different endings. The vocative always has the same form as the nominative. Fields, by colour: Some softening rules are (compared to the singular masculine nominative): See polski for sample "-ki" declension. Adjective forms can also get comparative and superlative forms. 

&lt; Hard and soft consonants &lt; ^ Polish ^ &gt; Masculine noun declension &gt; = Feminine noun declension = The declension of nouns in Polish is less regular than of adjectives, but follows a pattern that is in many ways similar to adjective declension. Let's take a look at a few typical feminine declension nouns in singular and are equal to nominative in plural): As you can see it's quite regular. Possible changes are: These changes aren't specific to the feminine noun declension - they happen throughout the Polish language, so you'd better get used to them. A bit less typical are feminine nouns that end in "-ia": Notice that in the genitive and the accusative the pronunciation is the same in the singular and plural. This is not usual in Polish, and may cause some problems if number is not obvious from context. One solution is to overemphasise difference between "e" and "ę" in speech (which are usually pronounced the same at word endings). A better solution is to use some adjective or pronoun, for example: And abstract feminine nouns ending in "-ść" (note vocative forms): There is some magic here with "ść" changing to "śc" but it's only spelling. You never write the softened version of a consonant before a vowel - you change it to the "normal" version and add "i" to mark it as "soft". You don't have to add "i" if it's already there. Without this magic endings would look like: Let's try to use that knowledge in practice. &lt; Hard and soft consonants &lt; ^ Polish ^ &gt; Masculine noun declension &gt; 

 ^ Polish ^ = Masculine noun declension = There are three masculine genders in Polish – masculine personal, masculine animate and masculine inanimate. The differences between them are regular: But there are also some other differences in declension. Before they get explained, let's take a look at some examples. Phonetic changes and irregularities are: Now some sentences:  ^ Polish ^ 

= Neuter noun declension = The neuter gender in Polish is least common. Many neuter words end in "-o" and their declension is somewhat similar to inanimate-masculine. The archaic instrumental oczyma doesn't have much to do with the neuter gender - it was the Indoeuropean dual number (for a pair of things), which disappeared in almost all modern Indo-European languages. You can still find it in some expressions. In the singular locative "krześle" you can see how softening works on a group of consonants – all consonants are softened, in that case "s" changes to "ś" and "ł" to "l". As with other declensions, "e" may appear between two final consonants if there is no vowel after them ("null" ending). The plural instrumental is "-ami"/"-iami" in almost all nouns, and "-imi"/"-ymi" in almost all adjectives. Here you can see one of the exceptions. For some historical reasons "a" disappeared and we have "-ćmi" instead of the expected "-ciami". In the case of "dziecko", "oko" and "ucho" you can see that a slightly different form is used as the base for the singular and plural forms. This is also quite archaic. If you find some neuter noun, try to decline it rather like "piwo", "krzesło" and "ciastko" rather than like "dziecko", "oko" or "ucho". Examples: Neuter nouns ending in "-e" or "-ę". The most common template is: Other templates are: 

= More on adjectives = Adjectives in Polish may be used in many ways, including: The English "X is Y" has many meanings. If Y is a noun, it usually means that X is a member of some class, for example "Kasia jest dziewczyną" (Kasia is a girl). In that case Y should be used in the instumentative. But if Y is an adjective, it means a completely different thing - that Y has a certain feature. For example "Kot jest czarny" (cat is black) doesn't mean that cat is member of class of black things. To mark that this is a different situation the nominative is used for Y. In the case of adjectives that can be used as nouns both cases may be used, but the nominative is much more common. Conjugation of the verb "to be": Adjectives used and their antonyms, in the feminine singular nominative Some examples with adjectives used as nouns: If an adjective describes some noun, it has to be used in the same case, number and gender as the noun. Here are some examples: You can also stack adjectives: 

We often want to group nouns or sentences, to say "X and Y" or "X or Y". In Polish the commonly used word for "and" is "i": The word "oraz" may also be used to express English "and": It sounds a little more sophisticated. However, you shouldn't use "oraz" inbetween adjectives - "Kot jest gruby oraz czarny" doesn't sound very natural. When in doubt, say "i". You can also use "z" + instrumental (with) and "bez" (without) + genitive When you group a few nouns, the result will always be plural, and if any of them was person-masculine, so will the whole expression be. There are many words for "or" - we mostly use "albo" or "lub" (the same meaning), but "czy" (which usually has a different meaning) may be used for this purpose too, especially in questions. ^ Polish ^ 

In Polish, sometimes yes/no questions are created by a rising tone, which is indicated by a question mark: However, this is more or less to the equivalent of saying in English: "Basia has a cat?" Another way to form a question is by using the word 'czy'. This makes it slightly more formal and also easier to understand for foreigners. Specific questions are created by question pronouns like: "kto" and "co" behave like nouns (but they don't have number) and "jaka" and "która" like adjectives, so you must use them in their proper form. A question pronoun is usually moved to the beginning of a sentence, even if it replaces something that would usually be at the end, like an object. ^ Polish ^ 

There are two classes of verbs in Polish - perfective verbs to talk about actions that are completed or will be completed, and imperfective verbs for actions that are taking place in some moment (no indication of completion). Very often they form pairs. Perfective verbs don't have a present tense - if something is already done you must use the past tense and if it's not - the future tense. This "future perfective" tense is identical in form to the present tense. Let's take the imperfective and perfective verbs for "to drink": To form other tenses besides the present and future perfective we need two more verb forms - the infinitive and the participle. The infinitive usually ends with "-ć". The participle depends on the subject's number and gender: Very often a perfective verb is formed from the imperfective by some prefix. Another class of pairs is verbs which are already derived from some other verb by prefixing. In Polish double prefixing is very rare, so an alternative form of base verb is used with the same prefix. Participles usually have regular endings showing the number and tense of the verb's subject: "ła", "ł", "ło", "li", "ły". In some cases "a" also changes to "e" in the plural person-masculine, where a softened consonant is used. This is one of the typical features of the Polish language, and it can also be found in other places, for example the locative of the noun "miasto" (city) is "mieście". The future imperfective is formed by the future form of the verb "to be" (equivalent to "will" in English) and either the third person form of the participle or simply the infinitive. The participial form is more common, however, if you're unsure of the correct ending, just use the infinitive. You can also use the verb "to be" on its own: Past tenses - past perfective and past imperfective - were originally formed in a similar way, but later the auxiliary verb merged with the participle. The Polish past perfective doesn't have anything to do with the English perfect tense - it's used for completed actions and is more similar to the English simple past. The participle alone is the correct third person past tense: In first and second person you have to add: Note that for the masculine singular ("-ł") you have to add an extra "-e-": There is also a tense to mark something that happened before a given event in the past, like the English "pluperfect", but it's very archaic now and rarely used in modern Polish. External links. ^ Polish ^ 

chcieć - to want lubić - to like umieć - can/to be able to potrzebować - to need You may use these verbs, and many others, together with other verbs in the infinitive: ^ Polish ^ 

Prepositions as hints to declensions. Prepositions, which generally show some sort of literal or figurative relationship between events, very rarely overlap exactly between different languages, both in literal and figurative usages. For example, to a non-native English speaker, it is not obvious why we "talk about" something, rather than talk "on", "at", or "in" something. The correspondence shown below between some key Polish and English prepositions should not be taken as a relation in terms of meaning (semantics); it is only a rough approximation. However, in Polish, both the choice of verb and the preposition (the Polish preposition, of course) are constrained by the declension. Once you develop some intuition as to which preposition you should use, or which declension, the relation between these can be used to guess the declension from the preposition or vice versa. Common links are as follows. Note that some prepositions, these regarding to place, can be used with two possible cases. Usually, a locative or instrumental are used when the acting takes place in this place and accusative when it is going towards the place. Compare: 

Indices.  __NOEDITSECTION__ 

Advanced. Verbs. grammar rules: 

Welcome to the Modern Greek WikiBook! This book is aimed at teaching from absolute beginners to advanced. It is a work in progress and is currently undergoing major restructuring - please see the discussion page or the restructuring page, please consider helping the project by contributing. 

&lt;includeonly&gt; Atomic Structure. &lt;/includeonly&gt; 

Particle Properties. Before learning about subatomic particles, some basic properties should be understood. Charge. Particles may be electrically charged. Charge is a property which defines the force that a particle will exert on other charged particles. There is a well known saying that applies perfectly: "Opposites attract." (Likewise, like charges repel.) Positive charges and negative charges will attract each other and come together. Two positive or two negative charges will push each other away. The amount of charge a particle has is measured in coulombs, but it is more conveniently expressed in terms of an integer. For instance, a helium ion that has 2 less electrons than usual has a charge of +2, and a bromide ion with one more electron than usual has a charge of -1. (This may seem backwards, but remember that an electron has a negative charge.) Notice that charge not only applies to subatomic particles, but also ions and other things as well. Always remember to specify if a charge is positive or negative. Unlike ordinary numbers, we "always write the plus sign for positive charges" to avoid confusion with a negative charge. Mass. Mass is the measure of inertia. From a subatomic point of view, mass can also be understood in terms of energy, but that does not concern us when dealing with chemistry. Mass for particles, atoms, and molecules is not measured in grams, as with ordinary substances. Instead, it is measured in "atomic mass units", or amu. For more information about mass and amu, read the previous chapters on properties of matter. The Nucleus. At the center of each atom lies the nucleus. It is incredibly small: if you were to take the average atom (itself miniscule in size) and expand it to the size of a football stadium, then the nucleus would be about the size of a marble. It is, however, astoundingly dense: despite the tiny percentage of the atom's volume it contains nearly all of the atom's mass. The nucleus almost never changes under normal conditions, remaining constant throughout chemical reactions. Nuclei are themselves made up of a pair of smaller and more dense particles, the proton and the neutron. These particles are collectively dubbed nucleons. Protons. Protons have a charge of +1 and a mass of 1 amu. They are often represented by a formula_1. Protons will be important when learning about acids and bases—they are the essence of acid. Remember that the number of protons in an atom is its "atomic number", and defines what element it will be. The number of protons in a nucleus ranges from one to over a hundred. Consider the element hydrogen. Its atomic number is 1, so it has one proton and one electron. If it is made into an ion (an atom with missing or extra electrons), it will simply be a lone proton. Thus, a proton "is" the nucleus of a hydrogen atom, and a proton "is" a hydrogen ion. Therefore, a proton can be written as formula_2 or formula_3, both symbols for a hydrogen ion. Neutrons. Neutrons have no charge and a mass of 1 amu. A neutron is slightly heavier than a proton, but the difference is insignificant. Neutrons are often written formula_4. Unlike the protons, neutrons cannot exist outside the nucleus indefinitely as they become unstable and break down. Within one nucleus there can be many protons and neutrons all in close proximity to one another. The number of neutrons in a nucleus ranges from zero to over a hundred. You may wonder why neutrons exist. They have no charge, so can they do anything? The answer is yes—neutrons are very important. Remember that opposites attract and likes repel. If so, then how can several protons stay clumped together in the dense nucleus of an atom? It would seem as if the protons would repel and scatter the nucleus. However, there is a strong nuclear force that holds the nucleus together. This incredible force causes nucleons to attract each other with much greater strength than the electric force can repel them, but only over extremely short distances. A delicate balance exists between the number of protons and neutrons. Protons, which are attracted to one another via the strong force but simultaneously repelled by their electromagnetic charges, cannot exist in great numbers within the nucleus without the stabilizing action of neutrons, which are attracted via the strong force but are not charged. Conversely, neutrons lend their inherent instability to the nucleus and too many will destabilize it. Lastly, neutrons are very important in nuclear reactions, such as those used in power plants. Neutrons act like a bullet that can split an atom's nucleus. Because they have no charge, neutrons are neither attracted nor repelled by atoms and ions. The Electron Cloud. Surrounding the dense nucleus is a cloud of electrons. Electrons have a charge of -1 and a mass of 0 amu. That does not mean they are massless. Electrons do have mass, but it is so small that it has no effect on the overall mass of an atom. An electron has approximately 1/1800 the mass of a proton or neutron. Electrons are written formula_5. Electrons orbit the outside of a nucleus, unaffected by the strong nuclear force. They define the chemical properties of an atom because virtually every chemical reaction deals with the interaction or exchange of the outer electrons of atoms and molecules. Electrons are attracted to the nucleus of an atom because they are negative and the nucleus (being made of protons and neutrons) is positive. Opposites attract. However, electrons don't fall into the nucleus. They orbit around it at specific distances because the electrons have a certain amount of energy. That energy prevents them from getting too close, as they must maintain a specific speed and distance. Changes in the energy levels of electrons cause different phenomena such as spectral lines, the color of substances, and the creation of ions (atoms with missing or extra electrons). Electron Interactions. "Atoms will always have equal numbers of protons and electrons", so their overall charge is zero. Atoms are neutral. Ions, on the other hand, are atoms that have gained or lost electrons and now have an unequal number of protons and electrons. If there are extra electrons, the ion will be negatively charged. If there are missing electrons, the ion will be positively charged, due to the majority of positive protons. Valence electrons (the outermost electrons) are responsible for an atom's behavior in chemical bonds. The "core electrons" are all of the electrons not in the outermost shell, and they rarely get involved. An atom will attempt to fill its valence shell. This occurs when an atom has eight valence electrons (as explained in the next chapter), so atoms will undergo chemical bonds to either share, give, or take the electrons it needs. Sodium, for example, is very likely to give up its one valence electron, so that its outer shell is empty (the shell underneath it is full). Chlorine is very likely to take an electron because it has seven and wants eight. When sodium and chlorine are mixed, they exchange electrons and create sodium chloride (table salt). As a result, both elements have full valence shells, and a very stable compound is formed. 

Electron shells. Each shell is subdivided into subshells, which are made up of orbitals, each of which has electrons with different angular momentum. Each orbital in a subshell has a characteristic shape, and is named by a letter. They are: s, p, d, and f. In a one-electron atom (e.g. H, He+, Li+2, etc.) the energy of each orbital within a particular shell is identical. However, when there are multiple electrons, they interact and split the orbitals into slightly different energies. The letters s, p, d and f specify the subshells (angular quantum number "l") and the orbital is specified by the magnetic quantum number "m." The angular and magnetic quantum numbers relate to the magnitude and direction of the electron's angular momentum, respectively. Within any particular subshell, the energy of the orbitals depends on the angular momentum of orbitals s, p, d, and f in order of lowest to highest energy. This image shows the orbitals (along with hybrid orbitals for bonding and a sample electron configuration, explained later). The s subshell. The simplest subshell in the atom is the 1s subshell. It has no radial or angular "nodes": the 1s subshell is simply a sphere of electron density. A node is a point where the electron positional probability is zero. As with all subshell the number of radial nodes increases with the principle quantum number (i.e. the 2s orbital has one radial node, the 3s has two etc.). Because the angular momentum quantum number is 0, there is only one choice for the magnetic quantum number - there is only one s orbital per shell. The s orbital can hold two electrons, as long as they have different spin quantum numbers. S orbitals are involved in bonding. formula_1 The p subshell. Starting from the 2nd shell, there is a set of p subshell. The angular momentum quantum number of the electrons confined to p subshell is 1, so each orbital has one angular node. There are 3 choices for the magnetic quantum number, which indicates 3 differently oriented p subshell. Finally, each orbital can accommodate two electrons (with opposite spins), giving the p orbitals a total capacity of 6 electrons. formula_2 The orbital of p subshell all have two lobes of electron density pointing along each of the axes. Each one is symmetrical along its axis. The notation for the orbitals of p subshell indicate which axis it points down, i.e. px points along the x axis, py on the y axis and pz up and down the z axis. Note that although pz corresponds to the "ml" = 0 orbital, px and py are actually mixtures of "ml" = -1 and "ml" = 1 orbitals. The p orbitals are degenerate — they all have the same energy. P orbitals are very often involved in bonding. The d subshell. The first set of d orbitals is the 3d set. The angular momentum quantum number is 2, so each orbital has two angular nodes. There are 5 choices for the magnetic quantum number, which gives rise to 5 different d orbitals. Each orbital can hold two electrons (with opposite spins), giving the d orbitals a total capacity of 10 electrons. formula_3 Note that all the d orbitals have four lobes of electron density, except for the dz2 orbital, which has two opposing lobes and a doughnut of electron density around the middle. The d orbitals can be further subdivided into two smaller sets. The dx2-y2 and dz2 all point directly along the x, y, and z axes. They form an eg set. On the other hand, the lobes of the dxy, dxz and dyz all line up in the quadrants, with no electron density on the axes. These three orbitals form the t2g set. In most cases, the d orbitals are degenerate, but sometimes they can split, with the eg and t2g subsets having different energy. Crystal Field Theory predicts and accounts for this. D orbitals are sometimes involved in bonding, especially in inorganic chemistry. The f subshell. The first set of f orbitals is the 4f subshell. There are 7 possible magnetic quantum numbers, so there are 7 f orbitals. Their shapes are fairly complicated, and they rarely come up when studying chemistry. There are 14 f electrons because each orbital can hold two electrons (with opposite spins). formula_4 

Filling Electron Shells. When an atom or ion receives electrons into its orbitals, the orbitals and shells fill up in a particular manner. Aufbau principle. You may consider an atom as being "built up" from a naked nucleus by gradually adding to it one electron after another, until all the electrons it will hold have been added. Much as one fills up a container with liquid from the bottom up, the orbitals of an atom are filled from the lowest energy orbitals to the highest energy orbitals. Orbitals with the lowest principal quantum number (formula_1) have the lowest energy and will fill up first, in smaller atoms. Larger atoms with more subshells will seem to fill "out of order", as the other factors influencing orbital energy become important. Within a shell, there may be several orbitals with the same principal quantum number. In that case, more specific rules must be applied. For example, the three p orbitals of a given shell all occur at the same energy level. So, how are they filled up? Answer: all the three p orbitals have same energy so while filling the p orbitals we can fill any one of the Px, Py or Pz first. It is a convention that we chose to fill Px first, then Py and then Pz for our simplicity. Hence you can opt for filling these three orbitals from right to left also. Aufbau principle state that “atomic orbitals are filled with electrons in order of increasing energy level”. Hund's Rule. According to Hund's rule, orbitals of the same energy are each filled with one electron before filling any with a second. Also, these first electrons have the same spin. This rule is sometimes called the "bus seating rule". As people load onto a bus, each person takes his or her own seat, sitting alone. Only after all the seats have been filled will people start doubling up. Pauli Exclusion principle. No two electrons can have all four quantum numbers the same. What this translates to in terms of our picture of orbitals is that each orbital can only hold two electrons, one "spin up" (+½) and one "spin down" (-½). Orbital Order. 1s, 2s, 2p, 3s, 3p, 4s, 3d, 4p, 5s, 4d, 5p, 6s, 4f, 5d, 6p, 7s, 5f, 6d, 7p, 8s. Although this looks confusing, there is an easy way to remember. Go in order of the lines from top to bottom, top right end to bottom left of each line. Understanding the above rules and diagrams will allow you to determine the electron configuration of almost any atom or ion. How to Write the Electron Configuration of an Isolated Atom. Electron-configuration notation is relatively straightforward. An isolated Calcium atom 20Ca, for example, would have configuration of 1s22s22p63s23p64s2 in its "ground state". Other configurations like 1s22s22p63s23p64s14p1 are possible, but these "excited states" have a higher energy. They are not stable and generally only exist for a brief moment. The ground state configuration for Ca could be abbreviated by using the preceding "noble gas" (the elements found all the way on the right of the periodic table) as [Ar]4s2, where Ar is argon. Noble gases have very stable configurations, and are extremely reluctant to lose or gain electrons. Noble gas atoms are also the only ones regularly found as isolated atoms in the ground state. Atoms of other elements all undergo bonding under the conditions that we live under and this affects the orbitals that the outermost electrons are in. In that sense the electron configurations for the other elements are somewhat hypothetical: to encounter an isolated atom of, say, tungsten (W), we would have to first vaporize a metal that boils at 5800K. However, knowing atomic configurations is useful because it does help us to understand how and why they bond, i.e. why and how they change the configuration of their outer "valence" electrons. Rule of Stability. A subshell is particularly stable if it is half full or full. Given two configurations, the atom would "choose" the more stable one. Example: In the following configuration, Cu: [Ar]4s23d9, copper's d shell is just one away from stability, and therefore, one electron from the s shell jumps into the d shell: [Ar]4s13d10. This way, the d shell is full, and is therefore stable, and the s shell is half full, and is also stable. Another example: Chromium has a configuration of [Ar]4s13d5, although you would expect to see four d electrons instead of five. This is because an s electron has jumped into the d orbital, giving the atom two half-full shells—much more stable than a d orbital with only four electrons. The stability rule applies to atoms in the same group as chromium and copper. If one of these atoms has been ionized, that is, it loses an electron, it will come from the s orbital rather than the d orbital. For instance, the configuration of Cu+ is [Ar]4s03d10. If more electrons are removed, they will come from the d orbital. Magnetism. Magnetism is a well-known effect. Chances are, you have magnets on your refrigerator. As you already know, only certain elements are magnetic. Electron configurations help to explain why. Diamagnetism is actually a very weak repulsion to magnetic fields. All elements have diamagnetism to some degree. It occurs when there are paired electrons. Paramagnetism is an attraction to external magnetic fields. It is also very weak. It occurs whenever there is an unpaired electron in an orbital. Both diamagnetism and paramagnetism are responses of spins acting "independently" from each other. This leads to rather weak repulsion and attraction respectively. However, when they are located in a solid they may also interact with each other and respond "collectively" and that can lead to rather different properties: Ferromagnetism is the permanent magnetism that we encounter in our daily lives. It occurs when all the unpaired spins in a solid couple and tend to align themselves in the same direction, leading to a strong attraction when exposed to a magnetic field. This only occurs at room temperature with three elements: iron (Fe), nickel (Ni), and cobalt (Co). Gadolinium (Gd) is a borderline case. It loses its ferromagnetism at 20oC; above that temperature the spins start to act alone. However, there are many alloys and compounds that exhibit strong ferromagnetic coupling. The strongest one is Nd2Fe14B Antiferromagnetism is also a permanent magnetism in which unpaired spins align, but they do so in opposite directions. The result that the material does not react very strongly to a magnetic field at all. Chromium (Cr) is an example. 'Ferrimagnetism is a combination of ferro- and antiferromagetism. Unpaired spins align partly in opposite directions, but the compensation is not complete. This is why the material is still attracted strongly to a magnetic field. Magnetite Fe3O4 is such a substance. It was the first material studied for its magnetic properties and may, well, be the one sitting on your fridge. 

The octet rule refers to the tendency of atoms to prefer to have eight electrons in the "valence shell" (outer orbital). When atoms have fewer than eight electrons, they tend to react and form more stable compounds. When discussing the octet rule, we do not consider d or f electrons. Only the s and p electrons are involved in the octet rule, making it useful for the "representative elements" (elements not in the transition metal or inner-transition metal blocks). An octet corresponds to an electron configuration ending with s2p6. Stability. Atoms will react to get in the most stable state possible. A complete octet is very stable because all orbitals will be full. Atoms with greater stability have less energy, so a reaction that increases the stability of the atoms will release energy in the form of heat or light. Reactions that decrease stability must absorb energy, getting hotter. The other tendency of atoms is to maintain a neutral charge. Only the noble gases (the elements on the right-most column of the periodic table) have zero charge with filled valence octets. All of the other elements have a charge when they have eight electrons all to themselves. The result of these two guiding principles is the explanation for much of the reactivity and bonding that is observed within atoms: atoms seek to share electrons in a way that minimizes charge while fulfilling an octet in the valence shell. Example. The formula for table salt is NaCl. It is the result of Na+ ions and Cl- ions bonding together. If sodium metal and chlorine gas mix under the right conditions, they will form salt. The sodium loses an electron, and the chlorine gains that electron. In the process, a great amount of light and heat is released. The resulting salt is mostly unreactive — it is stable. It won't undergo any explosive reactions, unlike the sodium and chlorine that it is made of. Why? Referring to the octet rule, atoms attempt to get a noble gas electron configuration, which is eight valence electrons. Sodium has one valence electron, so giving it up would result in the same electron configuration as neon. Chlorine has seven valence electrons, so if it takes one it will have eight (an octet). Chlorine has the electron configuration of argon when it gains an electron. The octet rule could have been satisfied if chlorine gave up all seven of its valence electrons and sodium took them. In that case, both would have the electron configurations of noble gasses, with a full valence shell. However, their charges would be much higher. It would be Na7- and Cl7+, which is much less stable than Na+ and Cl-. Atoms are more stable when they have no charge, or a small charge. Exceptions. There are few exceptions to the octet rule. Two Electrons. The main exception to the rule is hydrogen, which is at its lowest energy when it has two electrons in its valence shell. Helium (He) is similar in that it, too, only has room for two electrons in its only valence shell. Hydrogen and helium have only one electron shell. The first shell has only one s orbital and no p orbital, so it holds only two electrons. Therefore, these elements are most stable when they have two electrons. You will occasionally see hydrogen with no electrons, but H+ is much less stable than hydrogen with one or two electrons. Lithium, with three protons and electrons, is most stable when it gives up an electron. Less Than an Octet. Other notable exceptions are aluminum and boron, which can function well with six valence electrons. Consider BF3. The boron shares its three electrons with three fluorine atoms. The fluorine atoms follow the octet rule, but boron has only six electrons. Although atoms with less than an octet may be stable, they will usually attempt to form a fourth bond to get eight electrons. BF3 is stable, but it will form BF4- when possible. Most elements to the left of the carbon group have so few valence electrons that they are in the same situation as boron: they are "electron deficient". Electron deficient elements often show metallic rather than covalent bonding. More Than an Octet. In Period 3, the elements on the right side of the periodic table have empty d orbitals. The d orbitals may accept electrons, allowing elements like sulfur, chlorine, silicon and phosphorus to have more than an octet. Compounds such as PCl5 and SF6 can form. These compounds have 10 and 12 electrons around their central atoms, respectively. Even palladium has 18 electrons in its valance shell Odd Numbers. Some elements, notably nitrogen, have an odd number of electrons and will form somewhat stable compounds. Nitric oxide has the formula NO. No matter how electrons are shared between the nitrogen and oxygen atoms, there is no way for nitrogen to have an octet. It will have seven electrons instead. A molecule with an unpaired electron is called a "free radical" and radicals are highly reactive, so reactive that many of them only exist for a fraction of a second. As radicals go, NO and NO2 are actually remarkably stable. At low temperatures NO2 does react with itself to form N2O4, its "dimer", that is not a radical. 

Puzzles|Decision puzzles|Yet Another Weighing|Solution One way to solve this would be to try to manually work out all the possible combinations and their sums. That would take a while. Instead, notice that each item is half the weight of the previous item. Lets assume you have less items and work out the general rule for any number of items, where each is twice as much as the next item:  Items Weights Number of Combinations  1 5 1: 5  2 5, 10 3: 5, 10, 15  3 5, 10, 20 7: 5, 10, 15, 20, 25, 30, 35  4 5, 10, 20, 40 15: 5, 10, 15, 20, 25, 30, 35, 40, 45, 50, 55, 60, 65, 70, 75  5 5, 10, 20, 40, 80 Without even listing them, we can see the pattern. The numbers 1, 3, 7, 15 follow the formula of 2n - 1, where n=1, 2, 3, 4 i.e. the number of items. For 5 items, n=5, 25 - 1 = 31. The pattern would continue for more items:  6 - 63  7 - 127  8 - 255 

Before learning C syntax and programming constructs, it is important to learn the meaning of a few key terms that are central in understanding C. Block Structure, Statements, Whitespace, and Scope. Next we'll discuss the basic structure of a C program. If you're familiar with PASCAL, you may have heard it referred to as a block-structured language. C does not have complete block structure (and you'll find out why when you go over functions in detail) but it is still very important to understand what blocks are and how to use them. So what is in a block? Generally, a block consists of executable statements. But before we delve into blocks, let's examine statements. One way to describe statements is they are the text (and surrounding whitespace) the compiler will attempt to turn into executable instructions. A simpler definition is statements are bits of code that do things. For example: int i = 6; This declares an integer variable, which can be accessed with the identifier 'i', and initializes it to the value 6. The various data types are introduced in the chapter Variables. You might have noticed the semicolon at the end of the statement. Statements in C always end with a semicolon (;). Leaving off the semicolon is a common mistake many people make, beginners and experts alike! So until it becomes second nature, be sure to double check your statements! Since C is a "free-format" language, several statements can share a single line in the source file, like this: /* this declares the variables 'i', 'test', 'foo', and 'bar'  note that ONLY the variable named 'bar' is set to six! */ int i, test, foo, bar = 6; There are several kinds of statements. You've already seen some of them, such as the assignment (codice_1). A substantial portion of this book deals with statement construction. Back to our discussion of blocks. In C, blocks begin with an opening brace "{" and end with a closing brace "}". Blocks can contain other blocks which can contain their own blocks, and so on. Let's look at a block example. int main(void)  /* this is a 'block' */  int i = 5;  /* this is also a 'block', nested inside the outer block */  int i = 6;  return 0; You can use blocks with the preceding statements, such as the main function declaration (and other statements we've not yet covered), but you can also use blocks by themselves. Whitespace refers to the tab, space and newline characters that separate the text characters that make up the source code.&lt;br&gt; Like many things in life, it's hard to appreciate whitespace until it's gone. To a C compiler, the source code  printf("Hello world"); return 0; is the same as  printf("Hello world");  return 0; which is also the same as  printf (  "Hello world") ;  return 0; The compiler simply ignores most whitespace (except, for example, when it separates codice_2 from codice_3). However, it is common practice to use spaces (or tabs) to organize source code for human readability. Most of the time we do not want other functions or other programmer's routines accessing data we are currently manipulating, which is why it is important to understand the concept of scope. Scope describes the level at which a piece of data or a function is visible. There are two types of scope in C, local and global. When we speak of global scope, we're referring to something that can be seen or manipulated from anywhere in the program. Local scope applies to a program element that can be seen or manipulated only within the block in which it was declared. Let's look at some examples to get a better idea of scope. int i = 5; /* this is a 'global' variable, it can be accessed from anywhere in the program */ /* this is a function, all variables inside of it  are "local" to the function. */ int main(void)  int i = 6; /* 'i' now equals 6 */  printf("%d\n", i); /* prints a '6' to the screen, instead of the global variable of 'i', which is 5 */  return 0; That shows an example of local and global. But what about different scopes "inside" of functions?&lt;br&gt; /* the main function */ int main(void)  /* this is the beginning of a 'block', you read about those above */  int i = 6; /* this is the first variable of this 'block', 'i' */  /* this is a new 'block', and because it's a different block, it has its own scope */  /* this is also a variable called 'i', but in a different 'block',  because it's in a different 'block' than the first variable named 'i', it doesn't affect the first one! */  int i = 5;  printf("%d\n", i); /* prints a '5' onto the screen */  /* now we're back into the first block */  printf("%d\n", i); /* prints a '6' onto the screen */  return 0; Basics of Using Functions. Functions are a big part of programming. A function is a special kind of block that performs a well-defined task. If a function is well-designed, it can enable a programmer to perform a task without knowing anything about how the function works. The act of requesting a function to perform its task is called a function call. Many functions require a function call to hand it certain pieces of data needed to perform its task; these are called arguments. Many functions also return a value to the function call when they're finished; this is called a return value (the return value in the above program is 0). The things you need to know before calling a function are: Many functions use the return value for the result of a computation. Some functions use the return value to indicate whether they successfully completed their work. As you have seen in the intro exercise, the codice_4 function uses the return value to provide an exit status to the operating system. All code other than global data definitions and declarations needs to be a part of a function. Usually, you're free to call a function whenever you wish to. The only restriction is that every executable program needs to have one, and only one, main function, which is where the program begins executing. We will discuss functions in more detail in a later chapter, C Programming/Procedures and functions. The Standard Library. In 1983, when C was in the process of becoming standardized, the American National Standards Institute (ANSI) formed a committee to establish a standard specification of C known as "ANSI C". That standard specification created a basic set of functions common to each implementation of C, which is referred to as the Standard Library. The Standard Library provides functions for tasks such as input/output, string manipulation, mathematics, files, and memory allocation. The Standard Library does not provide functions that are dependent on specific hardware or operating systems, like graphics, sound, or networking. In the "Hello, World" program, a Standard Library function is used (codice_5) which outputs lines of text to the stream. 

Many Japanese verbs have pairs of "transitive" and "intransitive" verbs. In Japanese these are known as 他動詞 (other move verb) and 自動詞 (self move verb). Formally, the difference between these is that a transitive verb can take on a direct object, whereas an intransitive verb can not. There are a few pairs of distinct verbs in English that correlate to this: "raise"/"rise", "fell"/"fall" and "lay"/"lie". Transitive verbs can be thought of as more causitive which requires an agent to perform the action. Intransitive verbs move on their own and can be thought of as just existing. This is best explained by example. Contrast the following two sentences: The general patterns for transitive and intransitive sentences is: The topics of intransitive verbs are usually inanimate. Some pairings are listed in the following table: The rule of thumb is that intransitive verbs usually take nouns with the particles 「が」("ga") or 「は」("wa") that act as subjects, whereas transitive verbs take object nouns marked with 「を」("o"). Transitive verbs can also take a "ga"-subject or "wa"-subject, although it may be omitted. Note that some intransitive verbs can take an "o"-object that indicates a location. For example, 出る ("to leave") can be used with a direct object that is a location from which the subject is to leave from. See the table below for more examples: 

Glycolysis, literally meaning "to split sugar," is the initial step in any respiratory system. Glycolysis involves the breaking down of a sugar (generally glucose, although fructose and other sugars may be used) into more manageable compounds in order to produce energy. The net end products of glycolysis are two Pyruvate, two NADH, and two ATP (A special note on the "two" ATP later). Glycolysis is a process that all organisms undergo; and therefore the most fundamental and primitive of all energy production systems. Simplified Explanation. Glycolysis is a series of steps cells go through to transform sugar into energy that the cell can use. All cells are equipped to perform glycolysis, as it is the primary method cells make energy. The byproducts of glycolysis may be further digested to release more energy. glycosis occurs in the cytoplasm of a cell in a plant. glycosis is the first stage in Respiration. Energy Investment Phase. Glycolysis occurs in the cytoplasm of a cell after glucose is ingested through the process of phagocytosis, or "cell eating," in which the cell engulfs a solid compound. Once the glucose is inside the cytoplasm, one molecule of ATP splits and transfers a phosphate group to the glucose, or "phosphorylates" it, so that it becomes an ion which cannot leave the cell because the cell membrane is impermeable to ions. This step effectively allows cells to efficiently absorb and utilize glucose and keep any glucose molecules from escaping. The process of phosphorylation also makes the glucose chemically reactive. Once the glucose has been phosphorylated, it is called "Glucose-6-Phosphate," which will be abbreviated G6P. The G6P is rearranged into Fructose-6-Phosphate (F6P) by a protein. When the F6P is created, another ATP molecule splits and the F6P is phosphorylated by another phosphate group. This phosphate is attached to an opposite end of the F6P molecule which is now called Fructose-1, 6-Biphosphate. This is done so that the molecule will be not only more chemically reactive but so that it will split due to the tensions placed upon the molecule by the two phosphate groups. (Phosphate groups have a negative charge; this allows them to be added to create incredibly useful energy compounds, ATP among them. Because like charges repel each other, the two phosphate groups are placed in positions that allow the glucose molecule to split roughly in half, although an enzyme called aldolase is still used to perform the division.) The Fructose-1, 6-Biphosphate, aided by aldolase, is split into Dihydroxyacetone Phosphate (DP) and Glyceraldehyde-3-Phosphate (G3P), which are isomers of one another. The enzyme isomerase converts these molecules into one another. Because only G3P is used in the final stages of glycolysis, however, the reaction favors the conversion of DP into G3P. The overall effect of this reaction and the following steps is to send two molecules of G3P into the Energy Payoff phase of glycolysis. Energy Payoff Phase. The G3P molecules then attach to an enzyme that removes two electrons from each G3P, therefore "oxidizing" the molecules. (Please note that the two G3P molecules do not attach to the enzyme at the same time; rather, two of the same enzyme would be used to carry out this reaction.) One electron from each G3P molecule is transferred to a hydrogen ion (H+; a hydrogen ion is also a proton) and one electron from each G3P molecule is transferred to NAD+, a cellular cofactor whose purpose is to carry electrons. The addition of electrons to the H+ and NAD+ combines the two into NADH. Because there are two G3P molecules, there are 2 NADH produced in this step. Since energy is released from the G3P molecules to allow the for the synthesis of the two NADH, some of the remaining energy attaches an inorganic phosphate group in the cell to each G3P molecule, forming 1, 3 Biphosphoglycerate (1,3B). The two chemically reactive 1,3B molecules now lose the phosphate group that was just attached to them to an ADP, or Adenosine Diphosphate, molecule. The phosphorylation of ADP creates an ATP molecule; because there are two 1,3B molecules, two ATP are created. With this step, the initial ATP investment is repaid, and the 1,3B molecules become 3-Phosphoglycerate, or 3P. After this step is completed, another enzyme relocates the phosphate group on the 3P, changing the 3P into 2-Phosphoglycerate, or 2P. After the 2P is created, an enzyme called enolase "hydrolyzes" the 2P molecules, which means that one water molecule is removed from each 2P molecule, which means that 2 H2O have been produced because there are two 2P molecules. In this particular case, the 2P is hydrolyzed so that the resulting Phosphoenolpyruvate (PEP) is very unstable and reactive. The final step of glycolysis involves a protein called BS. Upon attaching to the Pyruvate Kinase proteins, the two PEP are divided into two Pyruvate and two additional phosphate group which are used to phosphorylate two ADP molecules into two ATP molecules. In prokaryotes, this all but concludes the respiration sequence. There is one final process that is undergone, however; fermentation. The process of glycolysis coupled with the process of fermentation is known as anaerobic respiration, "anaerobic" literally meaning "without air." There are two main types of prokaryotic fermentation: alcohol fermentation and lactate fermentation. Fermentation. In alcohol fermentation, the two pyruvate molecules lose two oxygen molecules and one carbon molecule are separated from the pyruvate and merge to become CO2 (Carbon Dioxide). The remaining molecule is called Acetaldehyde. The Acetaldehyde is "reduced" (or, gains electrons) by an NADH molecule. Because it is reduced, it increases in size. In this case, the NADH loses two electrons and two hydrogen ions, or protons, which returns it to its original NAD+ status. The two electrons and protons are then transferred to the Acetaldehyde, and it changes into Ethanol. Ethanol is the "alcoholic" ingredient in alcoholic beverages- in other words, a waste product of bacteria is one of the principal ingredients of beer, wine, and other drinks. In lactate fermentation, no carbon dioxide is released and the pyruvate is instead reduced by the oxidation of NADH into lactate. In more quantitative terms, two protons and two electrons are removed from two NADH and are instead attached to two pyruvate molecules to create two lactate molecules. This lactate is also known as lactic acid, which happens to cause muscle pain and fatigue in animals. One might wonder why the NADH would be reverted to its original NAD+ status, since energy was used to create it in the first place. The NADH is actually oxidized into NAD+ so that a "glycolytic cycle" can be preserved: if the NAD+ had not been replenished, the prokaryote would eventually run out of organic materials, and it would no longer synthesize ATP, which would mean its eventual death. In eukaryotes, glycolysis is only the beginning of respiration. Instead of undergoing fermentation, he products of glycolysis are sent into the Krebs Cycle, explained in the next section. 



The problems in the texts you have seen are for you to ensure that you understand the concepts and ideas explored. They are not intended to be very difficult, but understandably they are not very challenging. Questions here are intended for you to further use the ideas you have learnt to answer some more difficult questions. Some questions are relatively straightforward, some of these questions depend on different sections of this discrete mathematics text, some of these questions are meant to be examination-style questions. Do not be discouraged by the increase in difficulty - hints are sometimes available, and you will be able to increase your problem solving skills! Set theory questions. These questions depend on your knowledge of ../Set theory/. 

This annotated text both summarizes and analyzes Ayn Rand's epic novel, Atlas Shrugged. In addition, the Appendix singles out the significance of certain Characters, Companies, Concepts, Places, Technologies, and other things in the novel. Analyses and summaries. All quotes are taken from Ayn Rand. "Atlas Shrugged" (New York: Penguin 1992.). The author of this book is Ayn Rand. Other works by Ayn Rand include The Fountainhead, Anthem and We the Living. External links.  __NOEDITSECTION__ 

=Atlas Shrugged, Part 1= CHAPTER FIVE: THE CLIMAX OF THE D'ANCONIAS. Section 152: Part 1, Chapter 5, Section 2. In this long section, Dagny walks to the Wayne-Falkland Hotel to confront Francisco d'Anconia. While walking, she reminisces her childhood with Francisco and we learn why this man is so significant to her. 

=Atlas Shrugged, Part 1, Chapters 6-10= CHAPTER SEVEN: THE EXPLOITERS AND THE EXPLOITED. Section 172: Part 1, Chapter 7, Section 2. "Judgement to determine right and wrong, vision to see the truth, courage to act upon it, dedication to that which is good, integrity to stand by the good at any price." 

=Atlas Shrugged, Part 2, Chapters 1-5= CHAPTER FIVE: Account Overdrawn. Section 253: Part 2, Chapter 5, Section 3. Lillian meets Jim for lunch. He accuses her of double crossing him regarding Hank's behavior at his trial for illegally selling steel. She was supposed to control Hank. She tells him she didn't betray him, she just failed, and she doesn't yet know why. Lillian decides to meet Hank when his train comes in, but there is no car reserved for anyone named Hank Rearden. She realizes that Hank must be in a car under some other name, and suspects that he is having an affair. This pleases her, as she has been waiting for this moment. She saw that Hank was good and pure and hated him for it. She married him solely to destroy the good in him. 

=Atlas Shrugged, Part 2, Chapters 6-10= 

=Atlas Shrugged, Part 3, Chapters 1-5= CHAPTER ONE: Atlantis. Section 312: Part 3, Chapter 1, Section 2. . At the end of dinner, Dagny goes home with John Galt and spends the night at his house in “a room he never intended her to occupy”. 

=Atlas Shrugged, Part 3, Chapters 6-10= 

Atlas Shrugged is divided into three parts, with ten chapters in each part. Chapters are divided into any number of sections. 

Hello World - Writing, Compiling and Running a C++ Program. Below is an example of a simple C++ program: When you write a program, you use a development environment. Your development environment can be a basic text editor or a feature rich C++ integrated development environment (IDE). You should not use a word processor like Microsoft Word, because it adds formatting codes to the text. If a compiler is not already available to you, see the Where to get a compiler Section of the book. Open your development environment and type the program shown (or copy and paste it) and save it as hello.cc. Now compile it using the C++ compiler:  COMMAND_PROMPT&gt; "g++ hello.cc -o hello" The example uses GCC, the GNU Compiler Collection ( http://gcc.gnu.org ) but you could use any other compiler, or use an IDE to compile it. The above command would produce an executable called hello or hello.exe. Invoke the executable to run your first C++ program: "Unix:"  COMMAND_PROMPT&gt; "./hello"  Hello World!  COMMAND_PROMPT&gt; "Microsoft Windows:"  COMMAND_PROMPT&gt; "dear hello"  Hello World!  COMMAND_PROMPT&gt; Text that is "italicized" is typed by you and the bold text is output by the program. If you use an IDE, it might automatically color the code for you based on the syntax. Troubleshooting. You don't have the GNU C++ compiler installed. If you have a different compiler, check its documentation for the correct compilation command. Lot of weird errors, mentioning many times:  undefined reference to `std::basic_ostream' [..] Usually ending with:  collect2: ld returned 1 exit status To use g++ to compile your hello.cc, use:  "g++ hello.cc -o hello" For gcc, use:  "gcc hello.cc -o hello -lstdc++" You did not type the full path, try:  "./hello" Is there a hello program in this directory? Can you see it when you type ls? If not, your compilation (g++ hello.cc -o hello) failed or you have changed to a wrong directory. If you do not specify -o hello, g++ names the output file a.out (assembler output) for historical reasons. In such a case, type:  "./a.out" to execute the program. Your First C++ Program Explained. The preprocessing directive. Some features of C++ are part of the language and some others are part of a standard library. The standard library is a body of code that is available with every C++ compiler that is standards compliant. When the C++ compiler compiles your program it usually also links it with the standard C++ library. When you use features from the library, C++ requires you to "declare" the features you will be using. The first line in the program is a preprocessing directive. In our example it is shown bold and italicized:  #include &lt;iostream&gt; This line causes the C++ declarations which are in the iostream "header" to be included for use in your program. Usually the compiler inserts the contents of a header file called iostream into the program. Where it puts it depends on the system. The location of such files may be described in your compiler's documentation. A list of standard C++ header files is in the standard headers reference tables. The iostream header contains various declarations for input/output (I/O). It uses an abstraction of I/O mechanisms called streams. For example there is an output stream object called std::cout which is used to output text to the standard output. Usually, this displays the text on the computer screen. The "preprocessor" is a part of the compiler which does some transformations to your code before the "actual" compiler sees it. For example, on encountering a #include &lt;iostream&gt; directive, it replaces the directive with the contents of the iostream header file. main Function. int main()  // ... The lines above represent a block of C++ code, given the name main. Such a named block of code is called a function in C++ parlance. The contents of the block are called the "body" of the function. The word int is shown in bold because it is a keyword. C++ keywords have some special meaning and are also reserved words, i.e., cannot be used for any purpose other than what they are meant for. On the other hand main is not a keyword and you can use it in many places where a keyword cannot be used (though that is not recommended, as confusion could result). Every (standards-compliant) C++ program must define a function called main. This is where the execution of the program begins. As we shall see later, main may "call" other functions which may call yet other functions. The compiler arranges for main function to be called when the program begins executing. (Although this is generally true, it is not always true. There is an exception to main's being executed at the very beginning that we will see later.) Now let us look at the code inside the main function. Printing Hello World! The first line in main uses the std::cout object to print the "string" (sequence of characters) Hello World! and end the line:  std::cout « "Hello World!\n"; This line is a C++ statement. C++ statements are terminated by a semicolon (;). Within the statement «, called the "insertion operator" is used to output the string using the std::cout stream. C++ strings are enclosed within double quotes ("). The quotes themselves are not part of the string and hence not printed. The sequence \n is used within a string to indicate the end of the current line. Though the sequence is represented by two characters, it takes up only one character's worth of memory space. Hence the sequence \n is called the newline character. The actual procedure to start a new line is system-dependent but that is handled by the C++ standard library transparent to you. Modifications to the Above Program. Here is the same program with minor modifications: Comments. The line added at the beginning:  "// This program just displays a string and exits" is a comment that tries to explain what the code does. Comments are essential to any non-trivial program so a person who is reading the code can understand what it is expected to do. There is no restriction to what is contained between the comment delimiters. The compiler just ignores all that is there in the comment. Comments are shown italicized in our examples. C++ supports two forms of comments:  "/* This program displays a string"  " and then it exits */" Comments are also used at times to enclose code that we temporarily want the compiler to ignore, but intend to use later. This is useful in debugging, the process of finding out bugs, or errors in the program. If a program does not give the intended result, by "commenting out" code, it might be possible to track which particular statement has a bug. As C-style comments can stop before the end of the line, these can be used to "comment out" a small portion of code within a line in the program. Flushing the Output Stream Buffer. Whenever you "write" (i.e., send any output) to an output stream, it does not immediately get written. It is first stored in memory and may actually get written any time in the future. This process is called "buffering" and the regions in memory used for storing temporary data like this are called "buffers". It is at times desirable to "flush" the output stream buffers to ensure all data has been written. This is achieved by applying the insertion operator to an output stream and the object std::endl. This is what is done by the line:  std::cout « std::endl; Before flushing the buffer, std:endl also writes a newline character (which explains its name, end line). Hence the newline is omitted in the string printed in the previous line. Returning Success Code. In most operating systems, every program is allowed to communicate to the invoker whether it finished execution successfully using a value called the exit status. As a convention, an exit status of 0 stands for success and any other value indicates failure. Different values for the exit status could be used to indicate different types of failures. In our simple program, we would like to exit with status 0. C++ allows the main function to "return" an integer value, which is passed to the operating system as the exit status of the program. The statement:  return 0; makes main to return the value 0. Since the main function is required to return an "integer", the keyword int is used to begin the function definition. This statement is optional since the compiler automatically generates code to return 0 for the main function for the cases where control "falls off" without a return statement. This is why the first program worked without any return statements. Note that this is only a special case that applies only to the main function. For other functions you must return a value if they are declared to return anything. Common Programming Error 1 Though the return statement is optional, main should not be declared to return void (a function declared as "void" is a function which does not return anything) as in some other languages like Java. Some C++ compilers may not complain about this, but it is wrong. Doing this could be equivalent to returning just about any random number that happened to be stored in a particular memory location or register, depending on the platform. This practice can also be potentially damaging to some operating systems, which rely on the return code to determine how to handle a crash or other abnormal exit. Whitespace and Indentation. Spaces, tabs and newlines (line breaks) are usually called "whitespace". These are ignored by the compiler except within quotes, apart from the rule that preprocessing directives and C++-style comments end at a newline. So the above program could as well be written as follows: Note, however, that spaces are required to separate adjacent words and numbers. To make the program more readable, whitespace must be used appropriately. The conventions followed when using whitespace to improve the readability of code constitute an Indent style. For example, with alternate indent styles, the program could be written like this: or like this: 

Simple I/O. String manipulation. 1. Write a program that prompts the user for a string (pick a maximum length), and prints its reverse. 2. Write a program that prompts the user for a sentence (again, pick a maximum length), and prints each word on its own line. Loops. 1. Write a function that outputs a right isosceles triangle of height and width "n", so "n = 3" would look like 2. Write a function that outputs a sideways triangle of height "2n-1" and width "n", so the output for "n = 4" would be: 3. Write a function that outputs a right-side-up triangle of height "n" and width "2n-1"; the output for "n = 6" would be: Program Flow. 1. Build a program where control passes from main to four different functions with 4 calls. 2. Now make a while loop in main with the function calls inside it. Ask for input at the beginning of the loop. End the while loop if the user hits Q 3. Next add conditionals to call the functions when the user enters numbers, so 1 goes to function1, 2 goes to function 2, etc. 4. Have function 1 call function a, which calls function b, which calls function c 5. Draw out a diagram of program flow, with arrows to indicate where control goes Functions. 1. Write a function to check if an integer is negative; the declaration should look like bool is_positive(int i); 2. Write a function to raise a floating point number to an integer power, so for example to when you use it float a = raise_to_power(2, 3); //a gets 8 float b = raise_to_power(9, 2); //b gets 81 float raise_to_power(float f, int power); //make this your declaration Math. 1. Write a function to calculate if a number is prime. Return 1 if it is prime and 0 if it is not a prime. 2. Write a function to determine the number of prime numbers below n. 3. Write a function to find the square root by using . 4. Write functions to evaluate the trigonometric functions. 5. Try to write a random number generator. 6. Write a function to determine the prime number(s) between 2 and 100. Recursion. Merge sort. 1. Write a C program to generate a random integer array with a given length n , and sort it recursively using the Merge sort algorithm. - sorting a one element array is easy. - sorting two one-element arrays, requires the merge operation. The merge operation looks at two sorted arrays as lists, and compares the head of the list , and which ever head is smaller, this element is put on the sorted list and the head of that list is ticked off, so the next element becomes the head of that list. This is done until one of the lists is exhausted, and the other list is then copied onto the end of the sorted list. - the recursion occurs, because merging two one-element arrays produces one two-element sorted array, which can be merged with another two-element sorted array produced the same way. This produces a sorted 4 element array, and the same applies for another 4 element sorted array. - so the basic merge sort, is to check the size of list to be sorted, and if it is greater than one, divide the array into two, and call merge sort again on the two halves. After wards, merge the two halves in a temporary space of equal size, and then copy back the final sorted array onto the original array. Binary heaps. 2. Binary heaps : - given a position i (the parent) , i*2 is the left child, and i*2+1 is the right child. - ( C arrays begin at position 0, but 0 * 2 = 0, and 0 *2 + 1= 1, which is incorrect , so start the heap at position 1, or add 1 for parent-to-child calculations, and subtract 1 for child-to-parent calculations ). in descending order. Dijkstra's algorithm. Dijkstra's algorithm is a searching algorithm using a priority queue. It begins with inserting the start node with a priority value of 0. All other nodes are inserted with priority values of large N. Each node has an adjacency list of other nodes, a current distance to start node, and previous pointer to previous node used to calculate current node. Alternative to an adjacency list, is an adjacency matrix, which needs n x n boolean adjacencies. The algorithm basically iterates over the priority queue, removing the front node, examining the adjacent nodes, and updating with a distance equal to the sum of the front nodes distance for each adjacent node , and the distance given by the adjacency information for an adjacent node. After each node's update, the extra operation "update priority" is used on that node : "while" the node's distance is less than it's parents node ( for this priority queue, parents have lesser distances than the children), the node is swapped with the parent. After this, "while" the node is greater distance than one or more of its children, it is swapped with the least distant child, so the least distant child becomes parent of its greater distant sibling, and parent to the greater distant current node. With updating the priority, the adjacent node to the current node has a back pointer changed to the current node. The algorithm ends when the target node becomes the current node removed, and the path to the start node can be recorded in an array by following back pointers, and then doing something like a quick sort partition to reverse the order of the array , to give the shortest path to target node from the start node. Quick sort. 3. Write a C program to recursively sort using the Quick sort partition exchange algorithm. - an advantage of reuse is less writing time, debugging time, testing time. classes : those elements less than the partition value, the partition element, and everything above (and equal to ) the partition element. and one at the element next to the partition element , and repeatedly scanning the left pointer right, and the right pointer left. and these are swapped, which uses the temporary extra space. should occur before swapping, otherwise the right pointer may be swapping a less-than-partition element previously scanned by the left pointer. element just before the last element. This happens when all the elements are less than the partition. - if the right pointer is chosen to swap with the partition, then an incorrect state results where the last element of the left array becomes less than the partition element value. - if the left pointer is chosen to swap with the partition, then the left array will be less than the partition, and partition will have swapped with an element with value greater than the partition or the partition itself. - The left pointer stops on the first out of order element. The right pointer begins on the first out-of-order element, but the outer loop exits because this is the leftmost element. The partition element is then swapped with the left pointer's first element, and the two elements are now in order. - In the case of a 2 element in order array, the leftmost pointer skips the first element which is less than the partition, and stops on the partition. The right pointer begins on the first element and exits because it is the first position. The pointers have crossed so the outer loop exits. The partition swaps with itself, so the in-ordering is preserved. the plus one the intended initial scan positions, and use the pre-increment and pre-decrement operators e.g. ( ++i, --i) . 

Welcome to the Chinese wikibook, a free Chinese textbook on the Standard Mandarin dialect. This page links to lessons using simplified characters (used in mainland China, Singapore and Malaysia). There is also a Traditional Character Version available (used in Taiwan, Macau, and Hong Kong). Contributors.  __NOEDITSECTION__ 

=About Chinese= The Chinese language (汉语/漢語, 华语/華語 or 中文; Pinyin: Hanyu, Huayu, Zhongwen) is a member of the Sino-Tibetan family of languages. About one-fifth of the world speaks some form of Chinese as its native language, making it the most common natively-spoken language in the world. There is great internal variety within Chinese, and spoken Chinese languages such as Standard Mandarin (Putonghua), Shanghainese (Wu), and Cantonese, which are not mutually intelligible. Nevertheless, there is a single standardized form of Chinese known as Standard Mandarin, which is based on the dialect of Beijing, which is in turn its own Mandarin dialect, among a large and diverse group of Chinese dialects spoken in Northern and Southwestern China. Standard Mandarin is the official language of Mainland China and Taiwan, one of four official languages of Singapore, and one of six official languages of the United Nations. Standard Mandarin also corresponds to the modern standard written Chinese language used by people speaking all forms of Chinese from all corners of China, including Mandarin, Wu, Cantonese, Hakka, Min-nan, and so forth. This textbook will teach Standard Mandarin, both spoken and written. Chinese grammar is in many ways simpler than European languages (for example, you will see no tenses, plurals, or subject-verb agreement), but there are also plenty of pitfalls that will trip up the unsuspecting beginner (for example, you will encounter tones, measure words, and discourse particles, which do not feature as strongly in European languages). In addition, the complexity of the writing system often daunts newcomers, as Chinese is one of the few languages in the world that does not use an alphabet or a syllabary; instead, thousands of characters are used, each representing a word or a part of a word. However, most complex Chinese characters are composed of only a few hundred simpler characters and many contain phonetic hints. There is a common Western misconception of Chinese writing as having thousands of distinct and idiomatic symbols each representing a single word. However, Chinese writing is surprisingly mnemonic, granted it is not as simple as the writing of Romance languages. The government of China has developed a system of writing Standard Mandarin pronunciation in the Roman alphabet, known as "Hanyu Pinyin", or simply, "pinyin" (汉语拼音/漢語拼音, "spelling according to sounds"). Hanyu Pinyin is used to write out Chinese words phonetically in an effort to help learners of Chinese with their pronunciation. This wikibook will teach you Hanyu Pinyin first, before any actual sentences. All examples and new vocabulary will always be given together with Hanyu Pinyin. There are two character sets: Simplified Chinese characters (简体字/簡體字, Pinyin: Jiǎntǐzì) and Traditional Chinese characters (繁体字/繁體字, Pinyin: Fántǐzì). Traditional characters trace their lineage through thousands of years of Chinese history, and continue to be used in Hong Kong, Macau, Republic of China, and among many overseas Chinese. Simplified Chinese characters were the result of reforms carried out in Mainland China to increase literacy rates and is now used in Singapore as well. Many people may think that Simplified Chinese was made by the PRC government, but in fact many characters in Simplified Chinese were assembled from the calligraphy in ancient China. There is no denying however that some characters were made up recently. Two systems share many of the same characters or with systematic, predictable reductions in stroke; however, some changes are not as formulaic. As a result, most native Chinese speakers are able to write in only one of the two systems, though they can usually read both. You are recommended to do the same. It is considered easier for people who learn Traditional to read both sets than people who learn Simplified only, but Simplified characters are less intimidating for beginners. In this wikibook, all examples and vocabulary are given in both systems, and you are encouraged to choose one system and stick with it throughout. Chinese characters have also been used in the past by other neighbouring Asian countries, and are still being used by some of them today. Some older Koreans still know how to read and write Chinese characters, but although the members of younger generations are taught Chinese characters or "hanja", they are rarely used and unnecessary for literacy in Korean, with the native alphabet, "hangul". Chinese characters are occasionally used for abbreviations, to clarify technical vocabulary (as Chinese serves roughly the same role in Korean that Latin serves in English), and to write family and many personal names. The Japanese still preserve many Chinese characters or "kanji" today and use them along with two syllabaries to write the Japanese language. 

Sites in English (beginners). Wikibook's Very Own HSK Guide! WikiMedia Sites. See more at Wikipedia. 



Objective-C is an object-oriented programming language, and is a layer over the . This means that if you know how to write C, there are only a few syntax changes to learn. In this section, we will look at how we can implement classes and instantiate objects in Objective-C. If you are unfamiliar with object-oriented programming, see . Basic syntax. If you have studied C, you can skip this section and proceed to the next section, A word on runtimes. If not, or if you're a little rusty, read on. In C, code is contained within a function. Functions are composed of several statements, each of which are terminated by a semicolon. For example, a simple function in C to add two numbers may look like this: int add (int a, int b)  int c;  c = a + b;  return c; This sample has the following lines: For control flow, C and Objective-C use: The for, while, and do statements will continue execution of a loop until the condition is false. The switch and if statements jump to different statements depending on the condition. Objective-C does not implement or import any functions by default. Instead, they need to be imported by using the codice_1 preprocessor directive. (#include works too, but generally it is not used in Objective-C - #import works better because it won't import the same thing twice like #include can try to do) A word on runtimes. Objective-C requires what is known as a "runtime system" to provide you with Objective-C features. The runtime takes care of the creation, management, and destruction of objects. If you have the gcc compiler, you should have the Objective-C runtime already installed. Otherwise, you may have to install the runtime separately. Note that gcc is a compiler collection; the full install contains not only the C compiler, but also C++, Java, Objective-C and even Fortran-95 and Ada 2005 compilers. However, it is possible to install gcc only partially, for example only the C compiler. Therefore, the Objective-C runtime may or may not be installed with gcc. Check your distribution. Two main systems are used to run Objective-C programs: This text assumes you are using the GNU runtime, but the two runtime systems are almost identical. This text does not cover the OPENSTEP, Cocoa, or GNUstep frameworks, but the skills learned here will be helpful when developing with those systems. Writing classes in Objective-C. Writing an Objective-C class requires a few design decisions before we start writing any code. Say we are writing a class to represent a point called Point in a two-dimensional plane. We need to ask ourselves two questions: For this example, we'll use double variables for the "x" and "y" coordinates. We'll define a method to get both coordinates, and we'll define a method to get their distance from the origin. The interface. @interface Point : Object @private  double x;  double y; - (id) x: (double) x_value; - (double) x; - (id) y: (double) y_value; - (double) y; - (double) magnitude; @end Let's examine what each element of this interface means. @interface Point : Object The @interface line says that we begin the declaration of the Point interface. We inherit from another class called Object. Objective-C provides you with a generic class, called Object. The Object class is a "root class" -- it does not inherit from another class. The Object class provides a set of methods that provide key functionality for an object to be used and recognized by the Objective-C runtime. As in C, we need to include Object's header file, Object.h, before we can use the set of methods declared in the header. If you don't explicitly inherit from some class, your class becomes a root class. (In Objective-C, there can be many root classes.) That's probably not what you want to do, because creating a root class is a very tricky and advanced topic that is only useful in very specific situations. In most cases, you will want to inherit from Object or some class that inherits from Object, etc. If you develop in NeXTStep / GNUstep / Cocoa, you will mostly be using another root class called NSObject, which provides different basic methods than Object. The word import means that we only include the file "once". This solves problems like recursive includes. @private  double x;  double y; Anything between the braces in an interface declaration specifies the instance variables that the class has. The @private line is a visibility modifier: it says that the instance variables after it are "private", i.e. they are accessible only from the class that declares them. It is good practice to mark all your instance variables private: they contain the state of an object and they should never be changed except by the object itself. - (id) x: (double) x_value; The declarations for the methods come after the instance variables. This is a declaration for the method to set the x value. It's common Objective-C style to name the setter method with the same name as the variable it's setting. The hyphen specifies an instance method (we'll look at these later). Then comes the name for the method (this method is called x: - note the colon), and the colon signifies an argument, called x_value, and is of type double. The cast on the x_value argument is necessary, because unlike C, the default type is id, and not int. The type id is very special—it is a type that can hold "any" object whatsoever. The first cast is not strictly necessary, but should be used for clarity. The first cast tells us that the method x returns an object back. - (double) x; This is another method named x (no colon), but it returns a double. This is the specification for the method that gets the value of the x variable. There is no conflict with the previous method because the types are different, and the previous method takes one argument whilst this takes none. @end After all the methods and instance variables are specified, this symbol marks the end of the declaration. The interface specification goes in a .h file—a header file. It's customary to call the file after the class, so we would create a file called Point.h. The implementation. @implementation Point - (id) x: (double) x_value  x = x_value;  return self; - (double) x  return x; - (id) y: (double) y_value  y = y_value;  return self; - (double) y  return y; - (double) magnitude  return sqrt(x*x+y*y); @end This is the implementation for the Point class. We implement the methods in the interface defined above. Let's have a look at each element of the implementation in turn. Again, we import Point's interface, just as we do in C. @implementation Point This is a marker that identifies the beginning of the implementation. - (id) x: (double) x_value  x = x_value;  return self; This is a typical method implementation. We can use the x variable directly without having to declare it since it is already declared in the interface, and is accessible only to the methods of this class. The behaviour, in general, is like an ordinary C function. Here we assign the value of the argument x_value to the instance variable x. The function then, returns the entire, current object as modified. The keyword self represents the current object. - (double) x  return x; Here is the simple method to get the value of the x variable. We simply return it. The behaviour for the other methods should be similar to those above. The @end keyword ends the implementation. The implementation specification goes in a .m file. It's customary to call the file after the class, so we would create a file called Point.m. Using the objects. Since Objective-C comes from C, we write a main function to make use of the class that we just created. Here's one typical main function. int main(void)  Point *point = [Point new];  [point x:10.0];  [point y:12.0];  printf("The distance from the point (%g, %g) to the origin is %g.\n",  [point x], [point y], [point magnitude]);  return 0; Let's examine what happens here. We import the interface to Point so we can use it. We import stdio.h so we can use printf.  Point *point = [Point new]; This is a typical Objective-C method call:  [point x:10.0];  [point y:12.0]; These are some typical instance method calls. We call point's method x: and y:. In Objective-C terminology, we say that we send point a message to apply to the method x:. These messages assign point's x and y variables. Recall that the x: and y: methods had in them return self;. This means that we can chain the two messages together, as follows:  [[point x:10.0] y:12.0]; The message [point x:10.0] returns an object, point, with its x variable set. Then "on that object", the outer message assigns its y variable.  printf("The distance from the point (%g, %g) to the origin is %g.\n",  [point x], [point y], [point magnitude]); The printf statement has in it the method calls [point x], which returns the value of the x variable for printing, [point y], which does the same for the y variable for printing, and [point magnitude] which calculates the distance and returns that value. 

What is the maximum power that an appliance should have if it connect to a 230V supply by a 5A cable ? 

Well done this is the correct answer. On to the next question» 

Sorry this is an incorrect answer. Remember that to find the power you need to "multiply" the current and the voltage. «Go back and try again 

Sorry but this isn't the correct answer. Be careful about powers of ten when you when you do a calculation, especially when using a calculator. «Go back and try again 

Sorry but this is the wrong answer. Remember you need to "multiply" the voltage and the current together to get the power. «Go back and try again 

The graph above shows how the potential difference across a conductor varied with the current flowing through it. Which of the following statements is true ? 

Sorry this is the wrong answer.It is quite possible to answe this question from the information given. From Ohm's Law But V/I is also the "slope" (or "gradient") of the line on the graph. So with a constant slope, we must have a constant resistance. Go back and try again» 

All help is welcome. Trigonometry Book 1. Book 1 is pre-calculus trigonometry. We assume the student is relatively new to algebra and do algebra step by step. Many of the pages have closely related free/YouTube videos at the Khan Academy. This is by design. Many students find the video presentation helpful with learning mathematical material. As with all three trigonometry books, we have a "for Enthusiasts" section, which is for the student who finds the normal content and pace too slow and too easy, and yet still needs exercises and practice with Book 1 trigonometry. Trigonometry Book 2. Book 2 is also pre-calculus trigonometry. However, the algebra moves at a brisker pace than in Book 1. The topics are not central to understanding trigonometry as it is usually taught in schools, now that a lot of former content has been dropped. One rule of thumb of the topics in Book 2 is the union of the set of all topics in high-school contest related to trigonometry, applications, and the topics in the classical book "Plane and Spherical Trigonometry" by Palmer (link), subtracting any thoroughly discussed topics in Book 1, and excluding any topic that requires substantial use of calculus or the concept of limit (which should be done in Book 3). The topics are useful, for example, for students interested in maths contests. In the enthusiasts section there are topics and exercises that are useful to students who will go on to do work with computer graphics. Book 2 trigonometry deepens the understanding of the many relationships "between triangles and circles". It also shows how to tackle some harder trigonometric function identities. Where do These Belong? "This section is for Book 2 pages where we don't yet know how they should fit in. Teachers Notes. Scrap Heap. These are pages that are on the way out. Trigonometry Book 3. Book 3 uses and builds on calculus, complex numbers, matrices. We assume the student is relatively fluent with algebra. We will often combine simple steps to keep proofs/explanations short. Book 1 is a prerequisite, but book 2 isn't. There are many computing related topics, particularly in the "for Enthusiasts" section. Where do These Belong? "This section is for Book 3 pages where we don't yet know how they should fit in. Authors. , , , , Also thanks to the many contributors to mathematical articles on Wikipedia from which some of the content has been lifted. 

There are formula_1 in&lt;br&gt; a complete circle Units of Measure. We have been measuring angles in degrees, with formula_1 in a complete circle. However, what if we measured the circle according to how many units we went around it. Think about it this way, do you measure the runner going around the circular track according to the degrees from the centre or the meters around the circle? The obvious answer is meters around the circle. However, how do you measure this in trigonometry? Choice of Units for Length and Weight. In measuring many quantities we have a choice of units. For example with distances we can use the metric system and measure in metres, kilometres, centimetres, millimetres. It is also possible to measure distances in miles, yards, feet and inches. With weights we can measure in kilogrammes and grammes. We can also measure in pounds and ounces. Choice of Units for Measuring Time. In measuring time we choose to have 60 seconds in a minute and 60 minutes in an hour. We could devise a new more metric system for time and divide an hour into 100 units, each three fifths of our current minute, and then divide these shorter 'minutes' up into 100 units each of which would be about a third of a second. Why 60? Why 360? The choice of dividing into 60 is not entirely arbitrary. 60 can be divided evenly into 2,3,4,5 or 6 or 10 or 12 parts. 60 can't be divided evenly into 7 equal parts, each a whole number in size, but it's still pretty good. Using 360 degrees in a full circle gives us many ways to divide the circle evenly with a whole number of degrees. Nevertheless, we could divide the circle into other numbers of units. Metric Degrees? From the earlier talk of the metric system you might be anticipating that we are about to divide the circle up into 100 or 1000 'degrees'. There "is" actually a unit called the 'grade' or 'Gradian' (Grad on calculators which have it) in which angles are measured by dividing a right angle up into 100 equal parts, each of one Gradian in size. One Gradian is 0.9 of a degree - quite close to being one degree. The grade is in turn divided into 100 minutes and one minute into 100 seconds. This "centesimal" system (from the Latin "centum", 100) was introduced as part of the metric system after the French Revolution. The Gradian unit is nothing like as widely used as either degrees or the units that interests us most on this page. The unit we introduce here is called the Radian. Radians is the circumference measure at&lt;br&gt; the point from formula_3. Choice of Units for Radians. Radians are quite large compared to degrees (and to Gradians). There are about 6.28 Radians to a complete circle. There are about 57.3 degrees in one Radian. Are the statements: Compatible? It is not hard to check. We said "there are about 6.28 Radians to a complete circle". The exact number is formula_6, making the number of radians in a complete circle the same as the length of the circumference of a unit circle. Remember that: The circumference of a circle is&lt;br&gt; formula_7&lt;br&gt; where formula_8 is the radius. Justifying Choice of Units for Radians. At this stage in explaining trigonometry it is rather difficult to justify the use of these strange units. There aren't even an exact whole number of radians in a complete circle. In more advanced work, particularly when we use calculus they become the most natural units to use for angles with functions like formula_9 and formula_10. A flavour of that, but it is only a hint as to why it is a good unit to use, is that for very small angles. And the approximation is better the smaller the angle is. This "only works if we choose Radians as our unit of measure" and very small angles.  We claim that for small angles measured in radians the angle measure and the sine of the angle are very similar. Let us take one millionth of a circle. In degrees that is 0.00036 degrees. In Radians that is formula_12 Radians. The angle of course is the same. It's one millionth of a circle, however we choose to measure it. It is just as with weights where a weight is the same whether we measure it in kilogrammes or pounds. The sine of this angle, which is the same value whether we chose to measure the angle in degrees or in radians, it turns out, is about 0.00000628. If your calculator is set to use degrees then formula_13 will give you this answer. The Radian Measure. There are&lt;br&gt; formula_6 Radians&lt;br&gt; in a complete circle. It is traditional to measure angles in degrees; there are 360 degrees in a full revolution. In mathematically more advanced work we use a different unit, the radian. This makes no fundamental difference, any more than the laws of physics change if you measure lengths in metres rather than inches. In advanced work, If no unit is given on an angle measure, the angle is assumed to be in radians. formula_15 A notation used to make it really clear that an angle is being measured in radians is to write 'radians' or just 'rad' after the angle. Very very occasionally you might see a superscript c written above the angle in question. What You need to Know. For book one of trigonometry you need to know how to convert from degrees to radians and from radians to degrees. You also need to become familiar with frequently seen angles which you know in terms of degrees, such as formula_16 in terms of radians as well (it's formula_17 Radians). Angles in Radians are nearly always written in terms of multiples of Pi. You will also need to be familiar with switching your calculator between degrees and radians mode. Everything that is said about angles in degrees, such as that the angles in a triangle add up to 180 degrees has an equivalent in Radians. The angles in a triangle add up to formula_18 Radians. Defining a radian. A single radian is defined as the angle formed in the minor sector of a circle, where the minor arc length is the same as the radius of the circle. formula_19 Measuring an angle in radians. The size of an angle, in radians, is the length of the circle arc "s" divided by the circle radius "r". formula_20 We know the circumference of a circle to be equal to formula_21 , and it follows that a central angle of one full counterclockwise revolution gives an arc length (or circumference) of formula_22 . Thus formula_6 radians corresponds to formula_1 , that is, there are formula_6 radians in a circle. Converting between Radians and Degrees. Because there are formula_6 radians in a circle: To convert degrees to radians: formula_27 To convert radians to degrees: formula_28 Exercises.  

 Contributing. This is a wiki textbook -- feel free to edit it, update it, correct it, and otherwise increase its teaching potential. 

Original Author: Modified by 

The normal cell cycle consists of 2 major stages. The first is interphase, during which the cell lives and grows larger. The second is Mitotic Phase. Interphase is composed of three subphases. G1 phase (first gap), S phase (synthesis), and G2 phase (second gap). The interphase is the growth of the cell. The normal cell functions of creating proteins and organelles. The Mitotic Phase is composed of Mitosis and Cytokinesis. Mitosis, when the cell divides. Mitosis can be further divided into multiple phases. Cytokinesis, which is when the two daughter cells complete their separation. Mitosis is the division of the nucleus and cytokinesis is the division of the cytoplasm. There is some overlap between there two sub phases. Reproductive cell division is called meiosis, which yields a nonidentical daughter cells that have only one set of chromosomes. In other words, they have half as many chromosomes as the parent cell. Meiosis occurs in gonads, ovaries or testes. Therefore combining two gametes together produce 46 chromosomes. From Wikipedia. The cell cycle is the cycle of a biological cell, starting from the time it is first formed from a dividing parent cell until its own division into two cells, consisting of repeated mitotic cell division and interphase (the growth phase). A cell spends the overwhelming majority of its time in the interphase(about 90% of time). Background Information. DNA, deoxyribonucleic acid, consists of four nucleic acids, A, T, C, and G. In a cell, the DNA provides the directions for creating all of the proteins necessary for cell viability, health, growth, function, and replication. The unique DNA sequence that encodes each protein is called a gene, and the complete set of genes for an organism or cell is referred to as it's genome. Prokaryotic genomes are often a single long DNA molecule, and Eukaryotic genomes consist of number of DNA molecules. A typical human cell has about 2 m of DNA, which is 250,000 times greater than the cell's diameter. Before a cell divides the DNA is first copied then separated so that each daughter cell ends up with a complete genome. Chromosomes are the packaged DNA molecules. Because of chromosomes, the replication and distribution of so much DNA is manageable. Every eukaryotic species has a characteristic number of chromosomes in each cell nucleus. They contain two sets of each chromosome: one set inherited from each parent. For example human somatic cells (all body cells except the reproductive cells) each contain 46 chromosomes; the reproductive cells, gametes, have half as many chromosomes as somatic cells. The number of chromosomes in somatic cells varies widely among species. Eukaryotic chromosomes are made of chromatin that is a complex of DNA and associated protein molecules. Each single chromosome contains one very long, linear DNA molecule that carries several hundred to a few thousand genes; the associated proteins maintain the structure of the chromosome and help control the gene activity. When a cell is not dividing, each chromosome is a long thins chromatic fiber; however after DNA duplication chromosomes condense. Each chromatin fiber coils and folds. Each duplicated chromosome has two sister chromatids, containing an identical DNA molecule, initially attached along adhesive protein complex; such attachment is called sister chromatid cohesion. In condensed form of chromosome, a center narrow part is called centromere, a specialized region where the two chromatids are closely attached. The other part of a chromatid on either side of the centromere is referred as arm. Once the sister chromatids separate, they are considered individual chromosomes. Overview. The mitotic phase includes both mitosis and cytokinesis which is usually the shortest part of the cell cycle. Interphase accounts about 90%of the cycle; during interphase the cell grows and copies its chromosomes in preparation for cell division. Interphase is divided into sub-phases: G1 phase ("first gap"), the S phase ("synthesis"), and G2 phase ("second gap"). The chromosomes are duplicated only during the S phase. During G1 phase cell grows until S phase where the cell prepares for the cell division during G2 phase. Based on human cell, M phase only takes about 1 hour while the S phase occupies about 10-12 hours. The cell cycle consists of The cell cycle stops at several checkpoints and can only proceed if certain conditions are met, for example, if the cell has reached a certain diameter. Some cells, such as neurons, never divide once they become locked in a G0 phase. Mitosis. Mitosis has five stages: prophase, prometaphase, metaphase, anaphase, and telophase. Mitotic spindle starts to form in the cytoplasm during prophase. it is made of microtubules and other associated proteins. while the mitotic spindle assembles, the microtubules of the cytoskeleton disassemble, providing the material used to construct the spindle. In animal cells, the assembly of spindle microtubules starts at the centrosome, the microtubule-organizing center. In plant cells, the centrioles are not present. During interphase in animal cells, the single centrosome replicates; the two centrosomes remain together near the nucleus and they move apart during prophase and prometaphase of mitosis as spindle microtubules grows. The two centrosomes are located at the opposite end of the cell. Then aster, a radial array of short microtubules, extends from each centrosome. Kinetochore is a structure of proteins associated with specific sections of chromosomal DNA at the centromere. Each of the two sister chromatids of a replicated chromosome contains kinetochore as it face in opposite direction. During prometaphase, kinetochore microtubules form as come of the spindle microtubules attach to the kinetochores. After the microtubules are attached to chromosome's kinetochores, the chromosome begins to move towards the pole from which those microtubules extend. the chromosomes moves in a motion like a tug-of-war. Metaphase plate is the imaginary plane that formed during metaphase the centromeres of all the duplicated chromosomes are on the plane midway between the spindle's two poles. The other microtubules that did not attach to kinetochores overlap and interact with other nonkinetochore microtubules from the opposite pole. The nonkinetochore microtubules are responsible for elongating the whole cell during anaphase. During anaphase, the cohesins holding the sister chromatids of each chromosome are cleaved by enzymes. Then the chromatids separated, and they move towards the opposite ends of the cell. The region of overlap is reduced as motor proteins attached to the microtubules move away from one another, using ATP. As the microtubules push apart from each other, their spindle poles are pushed apart, elongating the cell. As the duplicate groups of chromosomes arrive at the opposite ends of the elongated parent cell, the telophase begins; during telophase nuclei reforms and cytokinesis begins. Cytokinesis. The cytokinesis process begins with cleavage. Cleavage furrow, a shallow groove in the cell surface near the old metaphase plate, is the first sign of cleavage. As it process, contractile ring of actin microfilaments form on the cytoplasmic side. The actin microfilaments interact with the myosin molecules, and cause the ring to contract. As the cleavage furrow deepens, the cell is separated into two with its own nucleus. For plant cells, there is no cleavage furrow because they have the cell walls. Instead of forming cleavages, vesicles derived from the Golgi apparatus move along microtubules to the middle of the cells, and forms cell plate. As the cell plate enlarges, and surrounding membrane fuses with the plasma membrane along the perimeter of the cell and from two daughter cells. Binary Fission. Binary fission is a method of asexual reproduction by "division in half". In prokaryotes, binary fission does not involve mitosis, but in single celled eukaryotes that undergo binary fission. In bacteria, motst genes are carried on a single bacterial chromosome that consists of a circular DNA molecule and associated proteins. The chromosome of the bacterium Escherichia coli, is 500 times as long as the cell when it is sctreched out. At the origin of replication, DNA of the bacterial chromosome begins to replicate. As the chromosome continues to replicate, one origin moves rapidly toward the opposite end of the cell, and the cell elongates. When the replication is complete the bacterium is about twice its initial size, and its plasma membrane grows inward, dividing the parent E. coli cell into two daughter cells. Bacteria don’t have mitotic spindles; the two origins of replication end up at opposite ends of the cell or in some other very specific location. The Evolution of Mitosis. Since the prokaryotes were on Earth more than a billion years than eukaryotes that mitosis had its origins in simpler prokaryotic mechanism of the cell reproduction can be assumed. Some of the proteins involved in bacterial binary fission are related to eukaryotic proteins that function in mitosis. Possible hypothesis of evolution of mitosis is that prokaryotic cell's reproduction gave rise to mitosis. The Cell Cycle Control System. Based from mammalian cell grow experiment, possible hypothesis was supported: the cell cycle is driven by specific signaling molecules present in the cytoplasm. In this experiment two cells in different phase of the cell cycle were fused to form a single cell with two nuclei. One cell was in the S phase and the other was in G1, and G1 nucleus immediately entered the S phase, as though stimulated by chemicals present in the cytoplasm of the first cell. Therefore, if a cell undergoing mitosis (M phase) was fused with another cell in any stage of its cell cycle, the second nucleus enteres mitosis. Other experiments on animal cells and yeasts demonstrates the sequential events of the cell cycle control system; the cell cycle control system operates set of molecules in the cell that both triggers and coordinates key events in the cell cycles. The cell cycle control system proceeds on its own, but it is regulated at certain checkpoints by internal and external signals. Animal cells have built-in stop signals that halt the cell cycle at checkpoints until they get go-ahead signals. The signals report whether crucial cellular processes that should have occurred by that point have in fact been completed correctly and thus whether or not the cell cycle should proceed. The three check points are in G1, G2, and M phase. For mammalian cells, G1 check points are the most important. When a cell receives a go-ahead signal at the G1 checkpoint, the cell complete the G1, S, G2 and M phases and divide; however when a cell does not get a go-ahead signal, it will exit the cycle and enter non dividing state, G0 phase. Most of human cells are in G0 phase, such as mature nerve cells and muscle cells. However the liver cells can re-enter the cycle by external signals such as growth factor released during injury. Rhythmic fluctuations in the abundance and activity of cell cycle control molecules pase the sequential events of the cell cycle. The regulatory molecules are portins of two types: protein kinases and cyclins. Portin kinases are enzymes that activate or inactivate other proteins by phosphorylating. The protein kinases give the go-ahead signals at the G1 and G2 checkpoints. The kinases that drive the cell cycle are present at a constant concentration in the growing cell, but they are in an inactive form. In order to activate them, kinase must be attached to a cyclin, a protein that cyclically fluctuating concentration in the cell. Because of such requirement, these are called cyclin-dependent kinases or Cdks. The activity of cdks rises and falls with changes in the concentration of its cyclin partner. The cylclin level rises during the S and G2 phases and then falls rapidly during M phase. MPF, the maturation -promoting factor, or M-phase -promoting factor, activity corresponds to the peaks of cyclin concentration. MPF triggers the cell's passage past the G2 checkpoint into M phase. MPF acts both directly as a kinase and indirectly by activating other kinases. During anaphase, MPF hels switch itself off by initiating a process that leads to the destruction of its own cyclin. The Cdk, noncyclin part of MPF, persists in the cell in inactive form until it associates with new cyclin molecules synthesized during the S and G2 phase of the next round of the cycle. Density-dependent inhibition is a phenomenon in which crowded cells stop dividing. It is caused by external physical factor. Also most animal cells exhibit anchorage dependence; in order to divide, the cells must be attached to a substratum; like a cell density, anchorage is signaled to the cell cycle control system via pathways involving plasma membrane proteins and elements of cytoskeleton linked to them. The loss of cell cycle controls leads to cancer cells, which exhibit neither density-dependent inhibition nor anchorage dependence. Reference. Berg, Jeremy M., John L. Tymoczko, and Lubert Stryer. Biochemistry. 7th ed. New York: W.H. Freeman, 2012. Print. Reece, Campbell, Lisa A. Urry, Michael L. Cain, Steven A. Wasserman, Peter V. Minosky, and Robert B. Jackson. Biology. 8th ed. San Francisco: Cummings, 2010. Print. 

Meiosis is a special type of cell division that is designed to produce gametes. Before meiosis occurs, the cell will be double diploid and have a pair of each chromosome, the same as before mitosis. Meiosis consists of 2 cell divisions, and results in four cells. The first division is when genetic crossover occurs and the traits on the chromosomes are shuffled. The cell will perform a normal prophase, then enter metaphase during which it begins the crossover, then proceed normally through anaphase and telophase. The first division produces two normal diploid cells, however the process is not complete. The cell will prepare for another division and enter a second prophase. During the second metaphase, the chromosome pairs are separated so that each new cell will get half the normal genes. The cell division will continue thorough anaphase and telophase, and the nuclei will reassemble. The result of the divisions will be 4 haploid gamete cells. Crossover. Crossover is the process by which two chromosomes paired up during prophase I of meiosis exchange a distal portion of their DNA. Crossover occurs when two chromosomes, normally two homologous instances of the same chromosome, break and connect to each other's ends. If they break at the same locus, this merely results in an exchange of genes. This is the normal way in which crossover occurs. If they break at different loci, the result is a duplication of genes on one chromosome and a deletion on the other. If they break on opposite sides of the centromere, this results in one chromosome being lost during cell division. Any pair of homologous chromosomes may be expected to cross over three or four times during meiosis. This aids evolution by increasing independent assortment, and reducing the genetic linkage between genes on the same chromosome. 

Mitosis is the normal type of cell division. Before the cells can divide, the chromosomes will have duplicated and the cell will have twice the normal set of genes. The first step of cell division is prophase, during which the nucleus dissolves and the chromosomes begin migration to the midline of the cell. (Some biology textbooks insert a phase called "prometaphase" at this point.)The second step, known as metaphase, occurs when all the chromosomes are aligned in pairs along the midline of the cell. As the cell enters anaphase, the chromatids, which form the chromosomes, will separate and drift toward opposite poles of the cell. As the separated chromatids, now termed chromosomes, reach the poles, the cell will enter telophase and nuclei will start to reform. The process of mitosis ends after the nuclei have reformed and the cell membrane begins to separate the cell into two daughter cells, during cytokinesis. The mitotic phase which includes both mitosis and cytokinesis is the shortest part of the cell cycle. The interphase cycle accounts for about 90% of the cell cycle. This phase is where the cell grows and copies its chromosomes in preparation for cell division. In the G1 phase which is also called the “first gap” the cell grows as it copies its chromosomes. In S phase, the cell starts to synthesize the DNA and completes preparation for cell division. In G2 it starts to divide. In biology, Mitosis is the process of chromosome segregation and nuclear division that follows replication of the genetic material in eukaryotic cells. This process assures that each daughter nucleus receives a complete copy of the organism's genetic material. In most eukaryotes, mitosis is accompanied with cell division or cytokinesis, but there are many exceptions, for instance among fungi. There is another process called meiosis, in which the daughter nuclei receive half the chromosomes of the parent, which is involved in gamete formation and other similar processes, which makes the parent cell still active. Mitosis is divided into several stages, with the remainder of the cell's growth cycle considered interphase. Properly speaking, a typical cell cycle involves a series of stages: G1, the first growth phase; S, where the genetic material is duplicated; G2, the second growth phase; and M, where the nucleus divides through mitosis. Mitosis is divided into prophase, prometaphase, metaphase, anaphase and telophase. The whole procedure is very similar among most eukaryotes, with only minor variations. As prokaryotes lack a nucleus and only have a single chromosome with no centromere, they cannot be properly said to undergo mitosis. Prophase. The genetic material (DNA), which normally exists in the form of chromatin condenses into a highly ordered structure called a chromosome. Since the genetic material has been duplicated, there are two identical copies of each chromosome in the cell. Identical chromosomes (called sister chromosomes) are attached to each other at a DNA element present on every chromosome called the centromere. When chromosomes are paired up and attached, each individual chromosome in the pair is called a chromatid, while the whole unit (confusingly) is called a chromosome. Just to be even more confusing, when the chromatids separate, they are no longer called chromatids, but are called chromosomes again. The task of mitosis is to assure that one copy of each sister chromatid - and only one copy - goes to each daughter cell after cell division. The other important piece of hardware in mitosis is the centriole, which serves as a sort of anchor. During prophase, the two centrioles - which replicate independently of mitosis - begin recruiting microtubules (which may be thought of as cellular ropes or poles) and forming a mitotic spindle between them. By increasing the length of the spindle (growing the microtubules), the centrioles push apart to opposite ends of the cell nucleus. It should be noted that many eukaryotes, for instance plants, lack centrioles although the basic process is still similar. Prometaphase. Some biology texts do not include this phase, considering it a part of prophase. In this phase, the nuclear membrane dissolves in some eukaryotes, reforming later once mitosis is complete. This is called open mitosis, found in most multicellular forms. Many protists undergo closed mitosis, in which the nuclear membrane persists throughout. Now kinetochores begin to form at the centromeres. This is a complex structure that may be thought of as an 'eyelet' for the microtubule 'rope' - it is the attaching point by which chromosomes may be secured. The kinetochore is an enormously complex structure that is not yet fully understood. Two kinetochores form on each chromosome - one for each chromatid. When the spindle grows to sufficient length, the microtubules begin searching for kinetochores to attach to. Metaphase. As microtubules find and attach to kinetochores, they begin to line up in the middle of the cell. Proper segragation requires that every kinetochore be attached to a microtubule before separation begins. It is thought that unattached kinetochores control this process by generating a signal - the mitotic spindle checkpoint - that tells the cell to wait before proceeding to anaphase. There are many theories as to how this is accomplished, some of them involving the generation of tension when both microtubules are attached to the kinetochore. When chromosomes are bivalently attached - when both kinetochores are attached to microtubules emanating from each centriole - they line up in the middle of the spindle, forming what is called the metaphase plate. This does not occur in every organism - in some cases chromosomes move back and forth between the centrioles randomly, only roughly lining up along the midline. Anaphase. Anaphase is the stage of meiosis or mitosis when chromosomes separate and move to opposite poles of the cell (opposite ends of the nuclear spindle). Centromeres are broken and chromatids rip apart. When every kinetochore is attached to a microtubule and the chromosomes have lined up along the middle of the spindle, the cell proceeds to anaphase. This is divided into two phases. First, the proteins that bind the sister chromatids together are cloven, allowing them to separate. They are pulled apart by the microtubules, towards the respective centrioles to which they are attached. Next, the spindle axis elongates, driving the centrioles (and the set of chromosomes to which they are attached) apart to opposite ends of the cell. These two stages are sometimes called 'early' and 'late' anaphase. At the end of anaphase, the cell has succeeded in separating identical copies of the genetic material into two distinct populations. Telophase. The nonkinetochore microtubules elongate the cell and try to cut the cell in two. The nuclear envelopes start to become created by fragments of the parents cell’s nuclear envelope. Then, the chromatids start to become less tightly coiled together. By this point, cytokinesis is fully under way. Cytokinesis. Cytokinesis refers to the physical division of one eukaryotic cell. Cytokinesis generally follows the replication of the cell's chromosomes, usually mitotically, but sometimes meiotically. Except for some special cases, the amount of cytoplasm in each daughter cell is the same. In animal cells, the cell membrane forms a cleavage furrow and pinches apart like a balloon. In plant cells, a cell plate forms, which becomes the new cell wall separating the daughters. Various patterns occur in other groups. In plant cells, cytokinesis is followed through by the usage of contracting ring of microfilaments that pull the cleavage furrow within itself, cutting the cell in two. In plant cells, vesicles from the Golgi apparatus start to form a cell plate within the center of the cell. When this cell plate solidifies and connects the two ends of the cell, a new cell wall is created and two daughter cells are produced. Regulation of Cell Cycle. Protein kinases are enzymes that activate or inactivate other proteins by phosphorylating them. These give out the signals for the G1 and G2 checkpoints to occur. However, to be active, the kinase must be attached to a cyclin. This is why it is called a CDK or a cyclin-dependent kinase. Internal kinetochores exhibit a wait function. Not until all kinetochores are attached to a spindle microtubule does the cell process starts. This helps prevent some chromosomes from being left behind. Density dependent inhibition is when cells have a cue to multiply until a certain level of density is fulfilled. This means that a cell keeps multiplying until there is a full layer or until a certain level of pressure is built upon each other. One possible explanation of why cancer cells do not follow normal signals is because they have an abnormality in the signaling pathway that conveys the growth factor’s signal to the cell-cycle control system. Usually, a cell will follow normal checkpoints due to the release of CDK in the system that regulate the cell process. However, in a cancer cell, the checkpoints are random. This means that because the cell does not follow density-dependent inhibition or follow the growth signals, the cell replicates at random points. 



About the platform. DOS, or Disk Operating System, can colloquially refer to any of a hundred different such operating systems. The name itself derives from it's ability to work with disks, a significant improvement over previous methods of storage. Generally this means the OS supplies a means of organizing, listing, reading, and writing files on the media. MS-DOS was Microsoft's first Operating System. It was built upon QDOS (Quick and Dirty Operating System), which was deliberated modeled after Gary Kildall's CP/M. The original MS-DOS was simplistic and very difficult to use for the untrained. As time progressed, the interface remained essentially the same (keyboard on a text console), however it had some significant usability features (e.g. DOSKey) implemented as well. MS-DOS retained the crown of most used DOS until Microsoft usurped its own OS with Windows, however there were other non-Microsoft disk operating systems (DOS) as well. DR-DOS was the primary competitor for MS-DOS and did fairly well until Windows 95 arrived. Afterwards, most people abandoned the non-GUI DOS system. A few diehards continued to use DOS, however, and some have produced a version of DOS written under an Open Source license known as FreeDOS. Although abandoned by its "creator", DOS is still a stable and viable operating system, although it has been overshadowed by Windows and the most popular open source operating system, Linux. DOS is still occasionally used on boot disks, so system recovery software may sometimes be written for DOS. The most popular languages for use on the DOS platform, besides DOS batch files and Intel x86 Assembly Language, are and . See QEMU/FreeDOS and A Neutral Look at Operating Systems/DOS for more information about FreeDOS. C/C++ Compilers. The main compiler for 32-bit DOS is DJGPP ( http://www.delorie.com/djgpp ). Nevertheless, 16 bit programs make up a substantial amount of the pre-Windows 95 program set, and so you may need to find a 16 bit compiler if you want to program on extremely old computers. For the most part however, one can just use a 32 bit DOS extender, such as CWSDPMI. Batch files. DOS allows the use of batch files, which are a collection or "batch" of DOS prompt commands stored in a file. When a user types in the filename (with or without the extension) of a batch file, DOS will perform each command listed in the .bat file, then return control to DOS. Assembly Language. DOS comes with a low-level debugger called DEBUG. This allows debugging of an executable program which it loads into memory along with DEBUG. This is done by running DEBUG "program" at the DOS prompt, where "program" is the name of your program. DEBUG will work with .com and .exe executables. You must include the file extension in the filename when you call it using DEBUG. DEBUG does not work with .bat files. DEBUG can also be run without a file to view CPU register contents, memory, and to assemble machine instructions directly into memory. A brief walkthrough of how to run DEBUG to view and change memory and registers, as well as assemble and run some basic machine instructions can be viewed here. 



The Russian Wikibook is a collaborative effort to create a comprehensive textbook for learners of the Russian language. Russian is an East Slavic language, related to Ukrainian and Belarusian, and is spoken by over 270 million people worldwide. This book includes four sections: a main text curriculum, a grammar supplement, an appendix, and a vocabulary. The main text guides the student through the lessons and provides everything to understand the texts that are to be understood. The grammar supplement provides a greater detail into the concepts presented in the lessons. The appendix is there to refer to for usage and other miscellaneous concepts. The vocabulary groups words into concept-based sections for studying. Содержа́ние (Contents). Слова́рь (Vocabulary). &lt;br&gt; &lt;br&gt; &lt;br&gt; &lt;br&gt; Miscellaneous resources Internet Resources Contributors 

These are some of the most commonly used words and phrases. Common phrases. ⁂ Click here for more information on Standard Indonesian phonology, how it differs from Standard Malay phonology. Vowel reduction is evident in modern Indonesian phonology due to the influence of Javanese. Helpful vocabulary for asking about words. When you're trying to find the right word, these can help you ask for suggestions: If you forget these, you can fall back on: "Sinonim" and "antonim" are not common words, but they are easy to remember, many Indonesian speakers understand them, and they can be helpful in your early stages of trying to communicate. Kata-kata Penting ("Important Phrases"). Lost traveller's guide for important phrases. Time. These are the words used to deal with almost any situation involving past or future. If neither word is used, it likely means that the action is happening around now. 

= Phrase = Grammar. Time and day. Comparisons. = Lessons = Miscellaneous. = Structure/Lesson Plans/Syllabus = Links and resources. These should be found place on pages where are relevant. Meta. Templates. Just an idea.... ... tell me what you think. Ruby. These are out-dated and should be replaced with when used in Japanese text, but kept for general use elsewhere on Wikibooks. Not sure what to do about these: 



ELEMENTS OF ART-INTRODUCTION. Audience (who the book is for). Students, producers and consumers of diverse forms of art, including but not limited to painting, drawing and sculpture, who would like to know more about the elements and theory of their specific art and art in general. This book is also for those who have an interest in the history of art. It will be richly illustrated, so that the facts and concepts it contains will be comprehended by children and adults, and be within their interest levels. Purpose (why it might be used). Elements of Art is designed to be used as a reference book. It is also a teaching book, to be consulted with other and more advanced/specialised resources on the subject, including textbooks that people may already be using. It can be a 'cram book' for examinations and tests, or a 'coffee table' book to take to the museum, library or keep at home. Form (the shape of the book). Chapters explaining (describing, analysing) the form, content and structure of particular elements of art, as well as general concepts reinforced throughout the book, and at the end of a section. Generally they are designed to build upon one another. Thus, there are chapters about line, colour, dimensionality (2-D forms of art like drawing and painting, and 3-D like sculpture) and shape. General Questions about Art (five Ws and a H-hopefully will be answered throughout the book, by clicking on specific links). WHO WHAT WHERE WHEN WHY and HOW What are the elements of art? How many elements? What are the forms of art? What is the difference between two-dimensional art? What is line? What is colour? How is line used? How is colour used? What is a dimension? What are the components of two dimensions? Of three? How can I show emotion? What have artists used to show emotion? What have artists used to show senses? Who is an artist? Who created the first art? Where was the first art created? Where were major art developments? Where can I go and see (hear/touch/taste/smell) art? Where can I go and learn more about art? Where do I start? When was the first art created? When were the major art developments? When is an artwork complete? Why do we create art? Why do we need art? Why do we respond to elements? Why use line/colour/whatever there and not here? http://SeaCloud9.org This web site is dedicated to the latest technological advancements, open source code, and cutting edge multimedia art. 

LINE. What is line? In art, line is the continuous movement of a mark from dot to dot. An identifiable path created by a point moving in space. Examples of line usage: Tughra of Sultan Suleiman the Magnificent Leonardo Da Vinci's The Head of the Virgin in Three-Quarter View Facing Right Francisco Goya's Plate 43 from 'Los Caprichos': The sleep of reason produces monsters Different lines. Lines of direction. The ability to manipulate a line includes suggesting its direction. There is no limit to a created line and below are the most common and basic types. Horizontal. Horizontal lines generally travel from left to right and are relative to the horizon. In art, it often establishes a feeling of rest as well as develops a ground in space. Vertical. Vertical lines travel up and down; they're perpendicular to horizontal lines. They often emphasize height and, in art, leads the eye from bottom to top and vice versa. Diagonal. Diagonal lines are angled and can either be an incline or decline slope. Artistically, they can be described as "unbalanced" and are considered neither horizontal nor vertical. Patterns can be created by using any of these lines, especially combining them, and manipulating them to produce variations. Lines of length. Lines can be short or long. Short lines. Short lines are lines that only extend for a short distance. Long lines. Long lines are lines that extend for a longer distance than a short line would. Lines of thickness. The thickness or thinness of a line can be achieved by using different materials. A ballpoint pen will create a thinner line than a big marker. Thick and thin lines are used to express different meanings in the work being creative. Using both Thick and thin lines together helps pull together a more cohesive work of art. Thick lines. Thick lines give the appearance of strength and allow a supportive quality to the lines around them. They tend to stand out and grab the eye's attention. Thin lines. Thin lines appear frail as if they can break under the slightest pressure. Thin lines give the piece a sense of elegance and lightness. Putting lines together. There are many ways to put lines together. One of the simplest is to join a horizontal and a vertical line. A square is two horizontal lines and two vertical lines. Many lines can be put together to make a shading effect. An example of this is cross-hatching, where many small horizontal lines are put over many small vertical lines. This gives a realistic effect. Putting curved lines together can create organic shapes unlike the vertical and horizontal lines which create a point. Effect of lines. Line creates movement and emotion in an artwork. Lines of movement. Lines can often give an artwork a sense of movement. For example, organic lines can create a sense of flowing movement, while geometric line can create a rigid feel or no movement in an artwork. Implied line can aid in guiding the viewers eye around an artwork, this is a form of movement in itself. The viewers eyes are moved by the implied line around the artwork. Repeated lines can create a vibration movement in an artwork. An example of implied line moving a viewer through an artwork can be seen in Raphael Sanzio's 'School of Athens'. The structure if the building create lines that unconsciously force the viewer to move through the painting and eventually to the focal point, the two men in the middle of the painting. As you can see there are 'lines' all leading to the direction of these men. Lines of emotion. There are a variety of different lines in art that all protrude some form of emotions. Organic, wavy lines create a mood of peacefulness and are softer on the viewer. While straight lines can have more harsh emotion for example power or anger. Vertical lines can also depict the same more hard emotions. Horizontal line are more peaceful than vertical lines. Diagonal lines create a sense of motion. An example of line protruding emotion can be seen in Van Gogh's 'Wheat field with Cypresses' The horizontal, curvy lines, created by the landscape and Van Gogh's visible brush strokes, sets the mood of the painting. The line creates a sense of calmness and tranquility. 

Chapter 15 Chapter 15. Magnoliophyta (I) The Division Magnoliophyta in the Kingdom Plantae comprises those species of plants that were formerly classified as angiosperms and are known widely as the "flowering plants". You have already studied flowers (Chapter 4), so now understand that the Division Magnoliophyta comprises all those species of plants that have flowers. For perspective: all of the plants we have read about in Section II of the "Botany Study Guide" up to this point do not have flowers, but certainly do have reproductive structures. We also know that flowers are the reproductive structures of the plants that bear them, and that reproductive structures are not limited to flowering plants. Thus, "flowers" are structures that distinguish plants in the Division Magnoliophyta from plants in the other divisions of the Kingdom Plantae. Observe a flower, and you know you are examining an angiosperm. However, not all angiosperms have obvious, showy flowers. You will need to consider that the structure of a flower is quite variable across all the many species of angiosperms (about a quarter million have been identified) and a few flowering plants actually seldom flower. However, angiosperm botanists put great stock in the structure of flowers as a way of classifying plants—more than any other part of a plant, the flower provides the basis for placement of a species in subtaxa (classes, orders, and families) of the Division Magnoliophyta. The Division Magnoliophyta is split into two large classes: the Magnoliopsida and the Liliopsida. There are over 300 families and 250,000 species of flowering plants. The remainder of this chapter will be concerned with the dicotyledonous angiosperms (Class Magnoliopsida), leaving the monocotyledonous angiosperms to be covered in the next chapter. Flower evolution. Conifers are pollinated by wind, meaning they must produce a large amount of pollen grains for only a few to arrive at the female megaspore, resulting in fertilization. An advancement that allowed for higher rates of fertilization per energy expended in making pollen would surely result in that new plant type being prolific and even dominant. However, a single mutation would never result in a flower adequately adapted to spread pollen using animals. Flowers are the result of a special kind of evolution called co-evolution. If plants reproduced and thrived with wind pollination, why did flowers evolve? Suppose one wind pollinated plant began to have its' pollen eaten by an animal, for example a beetle. Then suppose the beetle spreads the pollen from the male cones to the female cones while looking for pollen because it can't tell the difference between the male and female cones; it searches both, spreading pollen from the male to the female in doing so. This plant has an advantage over all wind pollinated plants, because of a higher rate of fertilization. Any mutation that leads to pollen being spread by an animal soon becomes prevalent within a species, because the plant pollinated by an animal is more efficient. The plant and the animal rely on each other, one for food, the other for reproduction. Any mutation in the animal that helps pollinate the plant will become prolific through natural selection, just as mutations that favor feeding the insect or making the insect come to the flower will be dominant. This back and forth evolution results in such seemingly improbable structures as flowers that look like female moths that are pollinated by amorous male moths, or flowers and bird beaks that complement the feeding of the bird and the pollination of the flowers. Class Magnoliopsida (dicots). Members of the Class Magnoliopsida are defined partly on the basis of the seed or seedling having two cotyledons, most obvious at germination. But the differences between dicots and monocots are many, and we will be able to recognize most flowering plants that we encounter as belonging to one or the other class without having to dissect the seed or observe the seedling. 

SHAPE. Shapes are created with lines in a given space, either real or imaginary. Shapes can be endlessly rotated. There can be organic shapes or geometric shapes. Different shapes. Circle. A circle is a shape with only one side created from a single, continuously curved line which encompasses the whole of the shape. Triangle. A triangle is a shape comprised of three straight lines which meet at three endpoints - the bottom side is horizontal, and the other two sides are diagonal, meeting each other at a point. Square. A square is a shape which is made of four straight lines which intersect at four points at 90 degree angles: the top and bottom lines are parallel to one another, as are the two lines comprising the sides of the square. In a square, each of the sides is the exact length of the other sides (a rectangle is a different shape where the opposing sides are equal in length; thus, all squares are rectangles but not all rectangles are squares.) Pentagon. A shape with 5 sides. The bottom side is horizontal, there are two vertical sides that are parallel and the two top sides are diagonal. A common use of the pentagon is to draw a house. Hexagon. A shape with six sides. 4 sides are diagonal and 2 are horizontal. Three-dimensional shapes. Three-dimensional shapes are not flat; instead, they create depth which creates form and the shapes appear touchable. Sphere. A round figure where every point on its surface an equal surface from the center. Examples: ball and globe. Cone. A solid or hollow object that tapers from a circular base to a point. It is a geometric shape formed from the base with lines that connect to one common point. Examples: funnel, traditional ice cream cone, traffic cones, and classic party hats. Cube. A cube has 8 endpoints, 12 edges and 6 faces. At every endpoint 3 lines intersect, and at an intersection any two edges are perpendicular to each other. Everything about the cube (edges, faces..etc) are equal. Think of a square with depth. Examples: Rubik cube and classic ice cube. Prism. A shape with two identical ends (often a shape) and flat sides that connect the ends. There are two common prisms: triangular and rectangular. However, a prism could be any shape as long as it is a polyhedron which means all faces are flat and all edges are straight This rules out a cylinder because it is curved. Putting shapes together. Most shapes in art are combinations of the shapes described above. They may be expressed (that is, they have a clear outline) or implied (the viewer has to seem them for his/herself). Also, different shapes can be put together for interesting results. Effects of shapes. Weight. Weight can capture the viewers eye by creating two-dimensional shape(s) that has a force that an element applies and attracts the eye. When it comes to shapes, you will notice that when you have shapes that are irregular, like an irregular triangle or quadrilateral, it will appear lighter than that of a regular shape. The reason for this being, is because irregular shapes make is appear as though part of the mass is taken away. When you put more elements into a space you are giving that space more weight. Height. Height has an effect on both two-dimensional forms and three-dimensional shapes. It related to how tall the shape can be made or stretch too. By having a variety of heights with shapes you are able to have all different types of proportions with how tall each shape is, one could be extremely tall while the other is shorter. 

Two and Three Dimensions. Dimensions, in art, can create some really interesting effects. They are connected with perspective, which is often used in graphic design, especially one-point perspective and two-point perspective. Two dimensions. The two dimensions are height and width. Three dimensions. The three dimensions are height and width and depth. Depth is often applied by projecting a shadow. Three-dimensional shapes can often be created by extra lines, or by doing a net. Art forms that use two dimensions. Drawing Painting Art forms that use three dimensions. Sculpture Origami Art forms that use two and three dimensions. Provide examples Mixed-media. 

The fundamental properties that we use to measure matter in are; Inertia, Mass, Weight, Volume, Density and Specific Gravity. The periodic table is a visual method of interpreting the chemical properties of elements which effect the measurements below. These measurements can be classified into two categories, intrinsic and extrinsic. the overall weight is equal to to another extrinsic properties Extrinsic properties (also called extensive), such as volume and weight, are directly related to the amount of material being measured. Density- the amount of how much an object/ matter is or how compact it is Intrinsic Properties. Intrinsic properties (also called intensive) are those which are independent of the quantity of matter present. For example, the density of gold is the same no matter how much gold you have to measure. Common intrinsic properties are density and specific gravity. 

... が好きです. ... ga suki desu Expresses that a thing or person is liked. The subject, marked by が, is liked by the person marked with the topic marker は (if mentioned). ... なければならない. ... nakereba naranai The construction "(verb stem)なければならない", "-nakereba naranai" is used to express an obligation: something that has to be done. Examples. Since ならない "naranai" is the plain-form negative of the verb なる "naru" (to become), in a more formal or polite situation one would use the polite negative form なりません "narimasen" instead. 下手の横好き（へたのよこずき）. heta no yoko-zuki An idiom meaning that you are "unskilled but enthusiastic" or "crazy about it, but not very good at it". Used about a hobby or skilled activity. 

Verb amar (to love) This is the typical example of a regular first conjugation verb in Spanish. 

Verb temer (to fear) This is the typical example of a regular second conjugation verb in Spanish. 

Verb partir (to leave "or" to cut) This is the typical example of a regular third conjugation verb in Spanish. 

Chapter 9 The term "algae" is used to collectively refer to a wide range (20,000-30,000 spp.) of very simple photosynthetic organisms. While this term is no longer used as a taxonomic grouping, it is still useful for referring informally to these photosynthetic protists. (Protists are diverse eukaryotes which are neither fungi, animals nor plants.) The earliest multicellular alga known is the red fossil alga Bangiomorpha (at right), found in 1,200 million year old rocks in Arctic Canada. What's more this is the first known organism to show sexual reproduction.== Chapter 9. Phycology ~ The Algae == The algae (singular: alga) comprise several different groups of plant-like organisms, some of which are (and some are not) regarded as members of the Kingdom Plantae. All algae lack true leaves, roots, flowers, and other structures found in the higher plants. They are distinguished from bacteria and protozoa mainly in that they are autotrophic, obtaining energy through photosynthesis. Although no longer considered a natural group, the term "algae" is still used for convenience. The botanical discipline concerned with the study of algae is called Phycology (or sometimes, Algology); and the environments most phycologists (or algologists) focus on are the marine intertidal/shallow subtidal regions of the world oceans. It is in these environments that the diversity of structurally complex algae (called seaweeds) reaches its pinnacle. As a grouping, the algae cut across even the prokaryote/eukaryote divide: the so-called "Blue-green algae" are cyanobacteria. All other algae are eukaryotes. Green Algae (different from Blue-green algae) are considered to be the ancestors of green plants. Other kinds of algae on the other hand are distinct from green plants and from each other in having different and unrelated accessory pigments. These pigments are responsible for the ways different algae absorb light, providing advantage to each individual type of alga to compete best at a water depth where its preferred wavelength is perhaps strongest. Cyanobacteria. The cyanobacteria comprise the structurally simplest algae, and presumably are closely related to the oldest photosynthetic organisms on the planet. Although capable of extracting energy from sunlight through photosynthesis, these algae are related to bacteria as evidenced by their prokaryotic cell structure. Yet, some Blue-greens have developed multi-cellular thalli that approach eukaryotic algal forms, and thus their traditional inclusion within the "algae." "Questions" 8.1 How do we know that the different algae are not monophyletic? 8.2 What do the different algae have in common that they are grouped together as algae? 8.3 What are 3 similarities between the green algae and green plants? 8.4 Define endosymbiosis. Did photosynthetic cyanobacteria exist before or after photosynthetic organelles? Why or why not? 

Primes and Modular Arithmetic. Factorisation Exercises. Factorise the following numbers. (note: I know you didn't have to, this is just for those who are curious) Recursive Factorisation Exercises. Factorise using recursion. Prime Sieve Exercises. 2. Find all primes below 200. 2 3 5 7 11 13 17 19 23 29 31 37 41 43 47 53 59 61 67 71 73 79 83 89 97 101 103 107 109 113 127 131 137 139 149 151 157 163 167 173 179 181 191 193 197 199 

HSE Primes|Primes and Modular Arithmetic. Factorisation Exercises. Factorise the following numbers. (note: I know you didn't have to, this is just for those who are curious) Recursive Factorisation Exercises. Factorise using recursion. Prime Sieve Exercises. 2. Find all primes below 200. 2 3 5 7 11 13 17 19 23 29 31 37 41 43 47 53 59 61 67 71 73 79 83 89 97 101 103 107 109 113 127 131 137 139 149 151 157 163 167 173 179 181 191 193 197 199 Division and Inverses Exercises. 1. 2. formula_30 3. 4. Coprime and greatest common divisor Exercises. 1. 2. We first calculate the gcd for all combinations Diophantine equation Exercises. 1. 2. 3. 4. Chinese remainder theorem exercises. 1. Question 1. Show that the divisible-by-3 theorem works for any 3 digits numbers (Hint: Express a 3 digit number as 100a + 10b + c, where a, b and c are ≥ 0 and &lt; 10) Solution 1 Any 3 digits integer "x" can be expressed as follows where a, b and c are positive integer between 0 and 9 inclusive. Now if and only if a + b + c = 3k for some k. But a, b and c are the digits of x. Question 2. "A number is divisible by 9 if and only if the sum of its digits is divisible by 9." True or false? Determine whether 89, 558, 51858, and 41857 are divisible by 9. Check your answers. Solution 2 The statement is true and can be proven as in question 1. Question 4. The prime sieve has been applied to the table of numbers above. Notice that every number situated directly below 2 and 5 are crossed out. Construct a rectangular grid of numbers running from 1 to 60 so that after the prime sieve has been performed on it, all numbers situated directly below 3 and 5 are crossed out. What is the width of the grid? Solution 4 The width of the grid should be 15 or a multiple of it. Question 6. Show that n - 1 has itself as an inverse modulo n. Solution 6 Alternatively Question 7. Show that 10 does not have an inverse modulo 15. Solution 7 Suppose 10 does have an inverse "x" mod 15, for some integer k but now "x" is not an integer, therefore 10 does not have an inverse 



Chapter 8 Chapter 8. Microbiology Microbes. Microbes are extremely small organisms or quasi-living particles traditionally studied by botanists, but now treated within the more specialized field of Microbiology. As a specialized field, Microbiology has its own methodologies and terminologies often very different from those used by botanists. Nonetheless, there is value for beginning botany students to learn something of these organisms—to gain understanding of their presumably more primitive nature and significance as the causes of plant diseases. Bacteria. Overview. A bacterium (plural: bacteria) is a single celled organism belonging to the Domain Bacteria, in the three domain scheme. It can also be a type of organism belonging to one of the three major branches of life. They belong to the kingdom monerans. Traditionally classified as one of the five kingdoms, bacteria are microscopic and relatively simple cells, called prokaryotes. These cells lack the nucleus and organelles of the more complex cells called eukaryotes that plants are constructed of. However, like the cells of plants, most bacteria possess a carbohydrate-based cell wall. In common speech, "bacteria" still refers also to archaeabacteria, although the latter recently have been classified as an independent branch or "domain" of life. They occur in various shapes which include: True bacteria are the oldest organisms on Earth, with the possible exception of the Archaea, and they are also the most abundant. Bacteria exist in soil, water, and as parasites of other organisms. Species and strains of bacteria cause many if not most non-hereditary diseases. They are the target of the drugs known as antibiotics. Viruses. Overview. A virus is an obligate cellular parasite that is completely dependent on the host cell for its replication. The genome of the virus may consist of single stranded or double stranded DNA or RNA. The size of viral genomes varies widely and may encode between one and 250 genes. Of most interest within the study of Botany are viruses that are plant pathogens. The majority of plant viruses have single-stranded, messenger-sense RNA genomes (Class IV) and encode only between one and 12 proteins. These proteins function in virus transmission, in replication, cell-to-cell and systemic movement, in the structure of the virus, and in the suppression of plant host defense mechanisms. In many cases, virus replication takes place in distinct virus-induced regions of the cell, the so-called viroplasms, and induces the synthesis of a pool of virus components followed by assembly of many virus particles from this pool. Viruses can usually be horizontally transmitted between hosts. In many cases the transmission is dependent on insects, nematodes, fungi or other vectors. However, some other viruses, for example "Tobacco mosaic virus" (TMV), are transmitted mechanically by physical contact between plant tissue and virus-contaminated surfaces. Once the virus is transmitted and has successfully entered into the plant cell, it moves locally from cell to cell until it enters the phloem, which allows the virus to enter distant tissue to cause a systemic infection. Even smaller than viruses are viroids . Viroids are infectious agents that consist of single-stranded RNA, which does not encode any proteins. History. The first virus to be identified was TMV. A. E. Mayer was a professor at the Agricultural College in Wageningen, the Netherlands that was established in 1876. Soon after he was asked by farmers to investigate a highly contagious disorder in tobacco, which he called "mosaic disease". In 1898, M. W. Beijerinck in the Netherlands concluded that the tobacco mosaic disease-causing agent was neither a bacterium nor any corpuscular body, but rather a contagium vivum fluidum, an infectious fluid. The next big step was made in 1935, when TMV was chemically purified by Nobel prizewinner Wendell M. Stanley in the United States. This was followed soon after in 1939 by the first visual observation of the rod shaped TMV particles by electron microscopy by Kausche, Pfankuch and Ruska in Germany. For a long time electron microscopy remained a main tool in virology, and many viruses were isolated and visualized. The study of the mechanisms of viral infection cycles became more important after molecular biology tools became available. Links to external portals / webpages. &lt;br&gt; 

To be able to study Trigonometry successfully, it is recommended that students complete Geometry and Algebra prior to digging in to the course material. Students should also be familiar with the arithmetic of the real number system. It is helpful to have a graphing calculator and graph paper on hand to be able to follow along as well. If one is not available, software available on sites such as GraphCalc or GeoGebra may be helpful. Geometric constructions proposed in the text can be drawn using Geops, free software for performing geometric constructions in the manner of the Ancient Greeks. "Next Page: In simple terms" 

Who do we mean when we speak of this person, Shakespeare? Shakespeare is William Shakespeare, one of the English-speaking world's greatest playwrights and poets, who possessed a great knowledge of human nature and transformed the English theatre. Yet many facts of his life remain a mystery. Some have been acquired from painstaking looks at the records of the time, so that this summary is based on generally agreed facts. It has been said that we only know three things about Shakespeare: that he was born, married and died. He was baptised on April 26, 1564; we do not know his birth date, but many scholars believe it was 23 April 1564. His father was John Shakespeare (who was a glover and leather merchant) and his mother Mary Arden (who was a landed local heiress). John had a remarkable run of success as a merchant, alderman, and high bailiff of Stratford, during William's early childhood. His fortunes declined, however, in the late 1570s. William lived for most of his early life in Stratford-upon-Avon. We do not know exactly when he went to London but he is said to have arrived in 1592. There is great conjecture about Shakespeare's childhood years, especially regarding his education. It is surmised by scholars that Shakespeare attended the free grammar school in Stratford, which at the time had a reputation to rival that of Eton. While there are no records extant to prove this claim, Shakespeare's knowledge of Latin and Classical Greek would tend to support this theory. In addition, Shakespeare's first biographer, Nicholas Rowe, wrote that John Shakespeare had placed William "for some time in a free school." John Shakespeare, as a Stratford official, would have been granted a waiver of tuition for his son. As the records do not exist, we do not know how long William attended the school, but certainly the literary quality of his works suggest a solid education. What is certain is that William Shakespeare never proceeded to university schooling, which has stirred some of the debate concerning the authorship of his works. In November 28, 1582, when he was 18, he married Anne Hathaway, who was 26. They had a daughter named Susanna, who was baptised on May 26, 1583. Later they had twins, a son named Hamnet and a daughter named Judith. Hamnet died while he was still a child on August 11, 1596. Due to the early death of his only son, Shakespeare does not have any direct descendants. For the seven years that followed the birth of his twins, William Shakespeare disappeared from all records, and then, turned up again in London some time in 1592. This period, which is known as the "Lost Years," has sparked as much controversy about Shakespeare's life as any period. When he was in London, he worked in repetory companies, and became part of the Lord Chamberlain's Men as an actor, playwright and shareholder. In 1599 he became an part-owner of the Globe Theater in Southwark. In 1603 James I became king and issued a royal licence to Shakespeare's acting company, who then became the King's Men, the foremost acting company in London at the time. In 1608 they leased a building called Blackfriars, which they converted to an indoor playhouse. It had some advantageous features like lighting and possibly heating. The Globe continued as their primary theater. From 1599 to 1608 he wrote several comedies and nearly all the famous tragedies. The year after (1609), his sonnets were published. William Shakespeare wrote his will in 1611, bequeathing his properties to his daughter Susanna (married in 1607 to Dr. John Hall). To his younger daughter Judith, he left £300, and to his wife Anne left "my second best bed." According to tradition, William Shakespeare died on his 52nd birthday, April 23, 1616. On his grave are the haunting words: "Good friend, for Jesus' sake forbeare" "To dig the dust enclosed here" "Blessed be the man that spares these stones," "And cursed be he that moves my bones." It took over 100 years for some of his bones to be stolen. After his death, in 1623, his friends published the First Folio, the first authorized collection of his works and a main source for the texts of his plays, 

 __NOEDITSECTION__ 

__NOEDITSECTION__ =Lesson 1: 你好！= It is appropriate to start off the introduction to Chinese with the common greeting 。 Below is a dialogue between two people meeting each other for the first time. Vocabulary. Note: Visit this lesson's Stroke Order subpage to see images and animations detailing how to write the following characters. Audio files of the words are linked from the pīnyīn when available. Problems listening? See . Proper Nouns. Forming the nationality is usually as simple as adding on to the country name. becomes , and so forth. Grammar. Basic Sentences. &lt;br&gt; 1. 我叫艾美。 Sentences using shì [是]. &lt;br&gt; &lt;br&gt; 1. 我是中国人。 2. 她是金妮。 3. 她们是英国人。 &lt;br&gt; &lt;br&gt; 1. 他不是东尼。 2. 我不是美国人。 Articles. There are no articles in Chinese grammar. While English noun clauses often begin with "a", "an", or "the", Chinese is less verbose. An example: An "a" appears in the English translation, but the singular and indefinite nature of is just inferred in Chinese. The question particle. The declarative example sentence in #1 is transformed into an interrogative in #2. 1. 她是金妮。 2. 她是金妮吗？ The question particle. 1. 我叫东尼, 你呢？ 2. 艾美是中国人, 他呢？ Question words. 1. 他们是哪国人？ 2. 谁是美国人？ 3. 她是谁？ 

Question 1. Is there a rule to determine whether a 3-digit number is divisible by 11? If yes, derive that rule.  Solution Let "x" be a 3-digit number We have now We can conclude a 3-digit number is divisible by 11 if and only if the sum of first and last digit minus the second is divisible by 11. Question 2. Show that "p", "p" + 2 and "p" + 4 cannot all be primes. ("p" a positive integer and is great than 3) Solution We look at the arithmetic mod 3, then "p" slotted into one of three categories Therefore "p", "p" + 2 and "p" + 4 cannot all be primes. Question 3. Find "x"  Solution Notice that Then Likewise, and Then Question 4. 9. Show that there are no integers "x" and "y" such that Solution Look at the equation mod 5, we have but therefore there does not exist "a" x such that Question 5. Let "p" be a prime number. Show that where E.g. 3! = 1×2×3 = 6 Hence, show that for "p" ≡ 1 (mod 4)  Solution a) If "p" = 2, then it's obvious. So we suppose "p" is an odd prime. Since "p" is prime, some deep thought will reveal that every distinct element multiplied by some other element will give 1. Since we can pair up the inverses (two numbers that multiply to give one), and (p - 1) has itself as an inverse, therefore it's the only element not "eliminated" as required. b) From part a) since "p" = 4"k" + 1 for some positive integer "k", (p - 1)! has 4"k" terms there are an even number of minuses on the right hand side, so it follows and finally we note that p = 4k + 1, we can conclude 

The romance category of plays contains four plays: They were all written late in Shakespeare's career. The Two Noble Kinsmen, of which Shakespeare was co-author, is sometimes included in this grouping. 

1.1: Points, Lines, Line Segments, and Rays Points and lines are two of the most fundamental concepts in Geometry, but they are also the most difficult to define. We can describe intuitively their characteristics, but there is no set definition for them: they, along with the plane, are the undefined terms of geometry. All other geometric definitions and concepts are built on the undefined ideas of the point, line and plane. Nevertheless, we shall try to define them. Point. A point is an exact location in space. A point is denoted by a dot. A point has no size. Line. As for a line segment, we specify a line with two endpoints. Starting with the corresponding line segment, we find other line segments that share at least two points with the original line segment. In this way we extend the original line segment indefinitely. The set of all possible line segments findable in this way constitutes a line. A line extends indefinitely in a single dimension. Its length, having no limit, is infinite. Like the line segments that constitute it, it has no width or height. You may specify a line by specifying any two points within the line. For any two points, only one line passes through both points. On the other hand, an unlimited number of lines pass through any single point. Ray. We construct a ray similarly to the way we constructed a line, but we extend the line segment beyond only one of the original two points. A ray extends indefinitely in one direction, but ends at a single point in the other direction. That point is called the end-point of the ray. Note that a line segment has two end-points, a ray one, and a line none. An angle can be formed when two rays meet at a common point. The rays are the sides of the angle. The point of the end of two rays is called the vertex. Plane. A point exists in zero dimensions. A line exists in one dimension, and we specify a line with two points. A plane exists in two dimensions. We specify a plane with three points. Any two of the points specify a line. All possible lines that pass through the third point and any point in the line make up a plane. In more obvious language, a plane is a flat surface that extends indefinitely in its two dimensions, length and width. A plane has no height. Space. Space exists in three dimensions. Space is made up of all possible planes, lines, and points. It extends indefinitely in all directions. N-dimensional Space. Mathematics can extend space beyond the three dimensions of length, width, and height. We then refer to "normal" space as 3-dimensional space. A 4-dimensional space consists of an infinite number of 3-dimensional spaces. Etc. Practice Problems. Conceptual Questions. Problem 1.1 (Geometry in Real Life) Give the geometric term(s) that is best modeled by each. a. The location of San Francisco, California&lt;br&gt; b. The surface of a chalkboard&lt;br&gt; c. The tip of a pencil&lt;br&gt; d. A piano chord&lt;br&gt; e. The edge of a desk&lt;br&gt; f. A knot in a rope&lt;br&gt; g. A telephone pole&lt;br&gt; h. Two connected walls&lt;br&gt; i. A partially opened folder&lt;br&gt; Problem 1.2 (Name the Plane) Use the appropriate notation to name the following plane in two different ways. Problem 1.3 (Name the Line) Use the appropriate notation to name the following line in five different ways. 

Lessons. This is a draft of the lesson plans. You can discuss these plans here. The second draft of this book is also being planned . Culture. "Văn hoá" 



 __NOEDITSECTION__ This is a work in progress, if you have suggestions or ideas to add in feel free but please do not completely rewrite what has already been written unless grossly inaccurate, if you have suggestions use the discussion page. Correct spelling, grammar, dates, names, places and any other such faux pas liberally; all spellings however are in Canadian English and not American. 

Here are some hints on how to express the date and time in English: 

The sound that "ng" and "ngh" make in Vietnamese is by far the hardest sound for Westerners to make. "Ng" and "ngh" simply make the last sound in " or " (as long as you don't make the hard /g/ sound at the end). The problem arises when "ng" or "ngh" come at the beginning of a word, as the common family name "Nguyễn" clearly demonstrates. Here, the speaker has to isolate the /ŋ/ sound, which even many Western dictionaries don't recognize in their pronunciation guides. (Those that do tend to represent it as /ng/.) This lesson will help you to at least pronounce the /ŋ/ sound well enough for a native listener. Oral exercise 1: Sing-along. For this exercise, follow the directions below. It may be a good idea to repeat each step until you've mastered it: Hopefully, you've just pronounced the word ""! It should sound something like . 

 you can used tamplets to help start your response paragraph. Ex: "One nation of the people, by the people, and for the people" is famous in part because Abraham Lincoln made a grammatically appealing statement. (Topic sentence) By using the prepositions ("of", "by", "for"), he reinforces his point by using parallel structure, and a repetitive vocative anaphora which harkens back to the role that the "people" play in making up a "nation". Reasons why this is a good passage: Vocative comes from the Latin (vocare) and is similar to other English word such as "vocation". It means "call". The "-ive" postfix normally means you are looking at a noun. Wikibookian Miss Winter AP English 11 September 2003 Free response: A person's identity is determined at birth 

Email spam (or just "Spam") is unsolicited email, similar to conventional "junk mail", but often on a larger scale. Spam is estimated to have cost U.S. organizations over $10 billion (in lost productivity) in 2003. Various software products are available to block or filter spam. Even lawmakers are stepping up to fight spam- in the United States the Can-Spam Act of 2003 was passed into law. "Spammers" - people or organizations responsible for sending out email spam - use large lists of email addresses that are collected in a number of ways. One of these means is by using computer programs that search websites for email addresses. An email address that appears on a website is therefore more likely to get spammed, other things being equal. There are various ways of preventing your email being harvested in this way. One is to include it as an image rather than a link. This is not ideal, though, because it means firstly that people will have to type the address themselves, and secondly, it makes the address inaccessible to blind people, or people who browse without images. An alternative, better method, is to use a small piece of 'JavaScript' to insert the email address into the page when it is displayed, keeping it out of the html which an email-collecting spam program might look at. One final method is to create a contact form to display on your website. The website user would fill out the form which when submitted forwards the message to you without displaying your email address. This final method has one drawback in that it circumvents the user's email system and does not provide the user with a record of the email that they sent. I have also seen many people try the format of me@REMOVETHISaol.com or whatever. The person who has a genuine interest in emailing you will remove the REMOVETHIS part of the eddress before sending, but programmes that gather email addresses en-masse will use it as-is - thus, you do not receive spam. One must be fairly inventive to make this effective; the instance given would be algorithmically handled by most site-scraping email bots. In addition, there is a new anti-spam feature now availible. Named the "Challenge/Response System", this either sends a link, or a word-verification page to a user, the first time they e-mail you. The user must either click the link or enter the word to verify they are not a spamming program. After this, you get the e-mail and they're added to your allow list. A study, by Brockmann &amp; Company IT consultants showed that challenge-response proved to be superior to appliances, hosted spam filters and commercial filters.Brockmann surveyed more than 500 businesses, with 40% of the respondents having IT responsibilities. The independently funded study resulted in the creation of a spam index to measure how satisfied workers were with their spam technologies. Despite being less sophisticated than filtering technology sold by antispam and antivirus vendors, the challenge-response method was twice as effective as hosted services for spam prevention. According to the survey, 67% of challenge-response users specified that they are very satisfied with their email experience as compared to next highest technology, hosted services, in which 42% reported that they were very satisfied. Commercial software filters, such as those produced by McAfee, Symantec and Trend Micro, scored the lowest satisfying only 22% of respondents.(SearchSecurity.com, Robert Westervelt Jul. 2007) www.bluebottle.com currently offers a public beta of this software. Spam is a very common thing and at one point or another it is something that we have all had. Spam is not only an annoying email, it is a tactics for marketing. However many times you have received this type of email, it is becoming more and more dangerous. It can be generated by businesses as well as individuals. It is used to promote products and is also brought about by things like forwarding ( things like jokes, images or chain letters) spam can be used to gather information from your computer if opened as well as used to send out viruses. It is becoming an increasing threat. Spam has become such an issue, that people now have to have an entire email address just as a "throw away account" an email address specifically catered toward junk mail because if we were to let it, it would take up about 90% of our inbox's, not only is this an issue for our email account but it is also being directed towards our phone, and needs to be handled with caution Email-spam-sample.png‎ Spam is a very common thing and at one point or another it is something that we have all had. Spam is not only an annoying email, it is a tactics for marketing. However many times you have received this type of email, it is becoming more and more dangerous. It can be generated by businesses as well as individuals. It is used to promote products and is also brought about by things like forwarding ( things like jokes, images or chain letters) spam can be used to gather information from your computer if opened as well as used to send out viruses. It is becoming an increasing threat. Spam has become such an issue, that people now have to have an entire email address just as a "throw away account" an email address specifically catered toward junk mail because if we were to let it, it would take up about 90% of our inbox's, not only is this an issue for our email account but it is also being directed towards our phone, and needs to be handled with caution 

Introduction. As soon as a child first learns about numbers, they become interested in big ones, a million, a billion, a trillion. They even make up their own, a zillion etc. One of the first mathematical questions a child asks is "what is the largest number?" This will often lead to a short explanation that there are infinitely many numbers. But there are many different "types" of infinity - in fact, there are infinite types of infinity! This chapter will try to explain what some of these types mean and the differences between them. Finite and Infinite Sets. There was once a mathematician called Georg Cantor who created a new branch of mathematics called set theory in the late 19th century. Set theory involves collections of numbers or objects. Here's a set: This set consists of five elements, namely the first five natural numbers. Now consider the set: Are these sets of the same size? Yes, they are. This is because they both have five elements. As we will see later, this method of comparing sizes does not work for all sets. An alternate method for comparing set sizes is to match elements of sets in a one-one fashion. Think of a small child who wants to compare the number of marbles she has with her brother's collection. Let's say she doesn't know how to count beyond ten. She can still compare the sizes of their collections of marbles by lining up their marbles in two parallel lines. The line on the left contains her marbles while the one on the right contains her brother's. If each marble on the left is aligned with exactly one marble on the right, then they both have the same number of marbles. We can use the same idea to compare infinite sets. If we can find a way to pair up one member of set A with one member of set B, and if there are no members of A without a partner in B and vice versa then we can say that set A and set B have the same number of members. Formally, two sets formula_3 and formula_4 are of the same size if there is a function formula_5 such that for every formula_6 in formula_3, we have formula_8 in formula_4 and moreover, for every formula_10 in formula_4, there exists an formula_6 in formula_3 such that formula_14. Example. Consider our previous example. We want to know if the sets formula_1 and formula_2 have the same size. We can create the following matching. Example. Let Set N be all counting numbers. N is called the set of natural numbers. 1,2,3,4,5,6... and so on. Let Set B be the negative numbers -1,-2,-3, ... and so on. Can the members of N and B be paired up? The formal way of saying this is "Can A and B be put into a one to one correspondence"? Obviously the answer is yes. 1 in set N corresponds with -1 in B. Likewise: and so on. Here, the one-one function that maps from A to B is formula_17. So useful is the set of counting numbers that any set that can be put into a one to one correspondence with it is said to be "countably infinite". Example. The set of integers is the set containing all elements from the set N, the set B and the element 0. That is The set of integers is usually denoted by Z. Note that N the set of natural numbers is a subset of Z. All members of N are in Z, but not all members of Z are in N. Is the set of integers countably infinite? In other words, can the set of integers be put in one-one correspondence with the set of all natural numbers? Since the set N is contained in the set Z, we may be tempted to declare that these two sets are not of the same size. However, we can and so on. We can write this one-one correspondence as a function formula_18 We can verify that this function generates all the integers in Z from the natural numbers in N. Strange indeed! A subset of Z (namely the natural numbers) has the same size as Z itself! Infinite sets are not like ordinary finite sets. In fact this is sometimes used as a definition of an infinite set. An infinite set is any set which can be put into a one to one correspondence with at least one of its subsets. Rather than saying "The number of members" of a set, people sometimes use the word cardinality or cardinal value. Z and N are said to have the same cardinality. Is the set of rational numbers bigger than N? In this section we will look to see if we can find a set that is bigger than the countable infinity we have looked at so far. To illustrate the idea we can imagine a story. There was once a criminal who went to prison. The prison was not a nice place so the poor criminal went to the prison master and pleaded to be let out. She replied: "Oh all right - I'm thinking of a number, every day you can have a go at guessing it. If you get it correct, you can leave." Now the question is - can the criminal get himself out of jail? (Think about if for a while before you read on) Obviously it depends on the number. If the prison master chooses a natural number, then the criminal guesses 1, the first day, 2,the second day and so on until he reaches the correct number. Likewise for the integers 0 on the first day, -1 on the second day. 1 on the third day and so on. If the number is very large then it may take a long time to get out of prison but get out he will. What the prison master needs to do is choose a set that is not countable in this way. Think of a number line. The integers are widely spaced out. There are plenty of numbers in between the integers 0 and 1 for example. So we need to look at "denser" sets. The first set that springs to most peoples mind are the fractions. There are an infinite number of fractions between 0 and 1 so surely there are more fractions than integers? Is it possible to count fractions? Let's think about that possibility for a while. If we try to use the approach of counting all the fractions between 0 &amp; 1 then go on to 1 - 2 and so on we will come unstuck because we will never finish counting the ones up to 1 ( there are an infinite number of them). But does this mean that they are uncountable ? Think of the situation with the integers. Ordering them ...-2, -1, 0, 1, 2, ... renders them impossible to count, but "reordering" them 0, -1, 1, -2, 2, ... allows them to be counted. There is in fact a way of ordering fractions to allow them to be counted. Before we go on to it let's revert to the normal mathematical language. Mathematicians use the term "rational number" to define what we have been calling fractions. A rational number is any number that can be written in the form p/q where p and q are integers. So 3/4 is rational, as is -1/2. The set of all rational numbers is usually called Q. Note that Z is a subset of Q because any integer can be divided by 1 to make it into a rational. E.g. the number 3 can be written in the form p/q as 3/1. Now as all the numbers in Q are defined by two numbers p and q it makes sense to write Q out in the form of a table. formula_19 Note that this table isn't an exact representation of Q. It only has the positive members of Q and has a number of multiple entries.( e.g. 1/1 and 2/2 are the same number) We shall call this set Q'. It is simple enough to see that if Q' is countable then so is Q. So how do we go about counting Q'? If we try counting the first row then the second and so on we will fail because the rows are infinite in length. Likewise if we try to count columns. But look at the diagonals. In one direction they are infinite ( e.g. 1/1, 2/2, 3/3, ...) but in the other direction they are finite. So this set is countable. We count them along the finite diagonals, 1/1, 1/2, 2/1, 1/3, 2/2, 3/1... Can we find any sets that are bigger than N? So far we have looked at N, Z, and Q and found them all to be the same size, even though N is a subset of Z which is a subset of Q. You might be beginning to think "Is that it? Are all infinities the same size?" In this section we will look at a set that is "bigger" than N. A set that "cannot" be put into a one to one correspondence with N, no matter how it is arranged. The set in question is R: the real numbers. A real number is any number on the number line. Remember that the set Q contains all the numbers that can be written in the form p/q with p and q integers, q different from 0. Most real numbers can never be put in this form and they are named "irrational numbers". Examples of irrational numbers include formula_21, formula_22 and formula_23. The set R is huge! Much bigger than Q. To get a feel for the different sizes of these two infinite sets consider the decimal expansions of a real number and a rational number. Rational numbers always either terminate: or repeat: Imagine measuring an object such as a book. If you use a ruler you might get 10cm. If you take a bit more care to and read the mm you might get 10.2cm. You'd then have to go on to more accurate measuring devices such as vernier micrometers and find that you get 10.235cm. Going onto a travelling microscope you may find its 10.235823cm and so on. In general the decimal expansion of any "real" measurement will be a list of digits that look completely random. Now imagine you measure a book and found it to be 10.101010101010cm. You'd be pretty surprised wouldn't you? But this is exactly the sort of result you would need to get if the book's length were rational. Rational numbers are dense (you find them all over the number line), infinite, yet much much rarer than real numbers. How we can prove that R is bigger than Q. It's good to get a feel for the size of infinities as in the previous section. But to be really sure we have to come up with some form of proof. In order to prove that R is bigger than Q we use a classic method. We assume that R is the same size as Q and come up with a contradiction. For the sake of clarity we will restrict our proof to the real numbers between 0 and 1. We will call this set R'. Clearly if we can prove that R' is bigger than Q then R must be bigger than Q also. If R' was the same size as Q it would mean that it is countable. This means that we would be able to write out some form of list of all the members of R' (This is what countable means, so far we have managed to write out all our infinite sets in the form of an infinitely long list). Let's consider this list. Where R1 is the first number in our list, R2 is the second, and so on. Note that we haven't said what order the list is to be written. For this proof we don't need to say what the order of the list needs to be, only that the members of R are listable (hence countable). Now lets write out the decimal expansion of each of the numbers in the list. Here r11 means the first digit after the decimal point of the first number in the list. So if our first number happened to be 0.36921... r11 would be 3, r12 would be 6 and so on. Remember that this list is meant to be complete. By that we mean that it contains "every" member of R'. What we are going to do in order to prove that R is not countable is to construct a number in R' that is not already on the list. Since the list is supposed to contain "every" member of R', this will cause a contradiction and therefore show that R' is unlistable. In order to construct this unlisted number we choose a decimal representation: Where a1 is the first digit after the point etc. We let a1 take any value from 0 - 9 inclusive "except" the digit r11. So if r11 = 3 then a1 can be 0, 1, 2, 4, 5, 6, 7, 8, or 9. Then we let a2 be any digit except r22 (the second digit of the second number on the list). Then a3 be any digit except r33 and so on. Now if this number, that we have just constructed "were" on the list somewhere then it would have to be equal to Rsomething. Let's see what Rsomething it might be equal to. It can't be equal to R1 because it has a different first digit (r11 and a1. Nor can it be equal to R2 because it has a different second digit, and so on. In fact it can't be equal to "any" number on the list because it differs by at least one digit from "all" of them. We have done what we set out to do. We have constructed a number that is in R' but is not on the list of all members of R'. This contradiction means that R' is bigger than any list. It is not listable. It is not countable. It is a bigger infinity than Q. Are there even bigger infinities? There are but they are difficult to describe. The set of all the possible combinations of any number of real numbers is a bigger infinity than R. However trying to imagine such a set is mind boggling. Let's look instead at a set that looks like it should be bigger than R but turns out not to be. Remember R', which we defined earlier on as the set of all numbers on the number line between 0 and 1. Let us now consider the set of all numbers in the plane from [0,0] to [1,1]. At first sight it would seem obvious that there must be more points on the whole plane than there are in a line. But in transfinite mathematics the "obvious" is not always true and proof is the only way to go. Cantor spent three years trying to prove it true but failed. His reason for failure was the best possible. It's false. Each point in this plane is specified by two numbers, the x coordinate and the y coordinate; x and y both belong to R. Lets consider one point in the line. 0.a1a2a3a4... Can you think of a way of using this one number to specify a point in the plane ? Likewise can you think of a way of combining the two numbers x= 0.x1x2x3x4... and y= 0.y1y2y3y4... to specify a point on the line? (think about it before you read on) One way is to do it is to take This defines a one to one correspondence between the points in the plane and the points in the line. (Actually, for the sharp amongst you, not quite one to one. Can you spot the problem and how to cure it?) Continuum hypothesis. We shall end the section on infinite sets by looking at the Continuum hypothesis. This hypothesis states that there are no infinities between the natural numbers and the real numbers. Cantor came up with a number system for transfinite numbers. He called the smallest infinity formula_24 with the next biggest one formula_25 and so on. It is easy to prove that the cardinality of N is formula_24 (Write any smaller infinity out as a list. Either the list terminates, in which case it's finite, or it goes on forever, in which case it's the same size as N) but is the cardinality of the reals = formula_25? To put it another way, the hypotheses states that: The hypothesis is interesting because it has been proved that "It is not possible to prove the hypothesis true or false, using the normal axioms of set theory" Further reading. If you want to learn more about set theory or infinite sets try one of the many interesting pages on our sister project . Limits "Infinity got rid of". The theory of infinite sets seems weird to us in the 21st century, but in Cantor's day it was downright unpalatable for most mathematicians. In those days the idea of infinity was too troublesome, they tried to avoid it wherever possible. Unfortunately the mathematical topic called analysis was found to be highly useful in mathematics, physics, engineering. It was far too useful a field to simply drop yet analysis relies on infinity or at least infinite processes. To get around this problem the idea of a "limit" was invented. Consider the series This series is called the harmonic series. Note that the terms of the series get smaller and smaller as you go further and further along the series. What happens if we let n become infinite? The term would become formula_29 But this doesn't make sense. (Mathematicians consider it sloppy to divide by infinity. Infinity is not a real number, you can't divide by it). A better way to think about it (The way you probably already do think about it, if you've ever considered the matter) is to take this approach: Infinity is very big, bigger than any number you care to think about. So let's let "n" become bigger and bigger and see if 1/"n" approaches some fixed number. In this case as "n" gets bigger and bigger 1/"n" gets smaller and smaller. So it is reasonable to say that the "limit" is 0. In mathematics we write this as and it reads: Note that we are not dividing 1 by infinity and getting the answer 0. We are letting the number "n" get bigger and bigger and so the reciprocal gets closer and closer to zero. Those 18th Century mathematicians loved this idea because it got rid of the pesky idea of "dividing by infinity". At all times "n" remains finite. Of course, no matter how huge "n" is, 1/"n" will not be "exactly" equal to zero, there is always a small difference. This difference (or error) is usually denoted by ε (epsilon). info -- infinitely small. When we talk about infinity, we think of it as something big. But there is also the infinitely small, denoted by ε (epsilon). This animal is closer to zero than any other number. Mathematicians also use the character ε to represent anything small. For example, the famous Hungarian mathematician Paul Erdos used to refer to small children as epsilons. Examples. Lets look at the function What is the limit as x approaches infinity ? This is where the idea of limits really come into its own. Just replacing "x" with infinity gives us very little: But by using limits we can solve it For our second example consider this limit as x approaches infinity of formula_34 Again lets look at the "wrong" way to do it. Substituting formula_35 into the expression gives formula_36. Note that you cannot say that these two infinities just cancel out to give the answer zero. Now lets look at doing it the "correct" way, using limits The last expression is two functions multiplied together. Both of these functions approach infinity as "x" approaches infinity, so the product is infinity also. This means that the "limit" does not exist, i.e. there is no finite number that the function approaches as "x" gets bigger and bigger. One more just to get you really familiar with how it works. Calculate: To make things very clear we shall rewrite it as Now to calculate this limit we need to look at the properties of sin(x). Sin(x)is a function that you should already be familiar with (or you soon will be) its value oscillates between 1 and -1 depending on x. This means that the absolute value of sin(x) (the value ignoring the plus or minus sign) is always less than or equal to 1: So we have 1/x which we already know goes to zero as x goes to infinity multiplied by sin(x) which always remains finite no matter how big x gets. This gives us Exercises. Evaluate the following limits; Infinite series. Consider the infinite sum 1/1 + 1/2 + 1/4 + 1/8 + 1/16 + ... Do you think that this sum will equal infinity once all the terms have been added ? Let's sum the first few terms. formula_46 Can you guess what formula_47 is ? Here is another way of looking at it. Imagine a point on a number line moving along as the sum progresses. In the first term the point jumps to the position 1. This is half way from 0 to 2. In the second stage the point jumps to position 1.5 - half way from 1 to 2. At each stage in the process (shown in a different colour on the diagram) the distance to 2 is halved. The point can get as close to the point 2 as you like. You just need to do the appropriate number of jumps, but the point will never actually reach 2 in a finite number of steps. We say that in the limit as n approaches infinity, Sn approaches 2. Zeno's Paradox. The ancient Greeks had a big problem with summing infinite series. A famous paradox from the philosopher Zeno is as follows: In the paradox of Achilles and the tortoise, we imagine the Greek hero Achilles in a footrace with the plodding reptile. Because he is so fast a runner, Achilles graciously allows the tortoise a head start of a hundred feet. If we suppose that each racer starts running at some constant speed (one very fast and one very slow), then after some finite time, Achilles will have run a hundred feet, bringing him to the tortoise's starting point. During this time, the tortoise has "run" a (much shorter) distance, say one foot. It will then take Achilles some further period of time to run that distance, during which the tortoise will advance farther; and then another period of time to reach this third point, while the tortoise moves ahead. Thus, whenever Achilles reaches somewhere the tortoise has been, he still has farther to go. Therefore, Zeno says, swift Achilles can never overtake the tortoise. 



Welcome to the Korean Wikibook, a free textbook for learning Korean. Note: To use this book, your web browser must first be configured to display Korean ("Hangeul") characters. Check the two boxes below: The boxes show "Hangeul" characters and "jamo". If symbols appear as blank boxes, garbage, or question marks (?), your computer or web browser needs to be configured for the Korean language. Introduction. Korean is the official language of both the Democratic People's Republic of Korea (North Korea) and the Republic of Korea (South Korea). It is also one of the two official languages in the Yanbian Korean Autonomous Prefecture. Worldwide, there are about 80 million Korean speakers, most of which (outside Korea) live in China, Japan or the United States, but they also represent sizeable minorities in Russia (esp. Far East and Sakhalinsk), New Zealand, Kazakhstan, Canada, Uzbekistan and Australia. In the Republic of Korea, the language is most often called 한국말 ("Han-guk-mal"), or more formally, 한국어 ("Han-guk-eo") or 국어 ("Guk-eo"; literally "national language"). In North Korea and Yanbian, the language is most often called 조선말 ("Chosŏnmal"), or more formally, 조선어 ("Chosŏnŏ"). Experts are still not completely sure of the origins of the Korean language, although it is generally believed to come from the Altaic language tree. It is an agglutinative language, so it has some certain special characteristics that are unlike English. A student of Chinese languages will quickly notice that Korean shares much of their vocabulary, while a Japanese student will also notice similarities in grammar and vocabulary. Feel free to use English Wiktionary's Korean language Category as a reference for these courses. New students to this type of language may initially progress slowly, but as study progresses, previously unfamiliar aspects of Korean will begin to make sense and new concepts will be more easily learned. Korean grammar is complex but surprisingly also very simple, and always very fun to learn. 



../Verbs/ =Ser= "Ser" is an irregular verb that means "to be". Indicative. ../Future/. — "or" — Subjunctive. ../Imperfect Subjunctive/. — "or" — Pluperfect subjunctive. — "or" — Compound tenses. These tenses are combinations of the above tenses. 

../Verbs/ No bones about it. The present tense is EASY. Just memorize the endings for the different verbs and a couple of oddball exceptions and you will be able to talk about any momentary thing that is on your mind. Who needs to plan for the future, dwell on the past, or delve into abstraction anyway ? There are three verb types in Spanish. They are called -AR, -ER, and -IR verbs. Each of these refers to the last two letters of the infinitive form of the verb. Examples: -AR: bailar, "to dance"; -ER: aprender, "to learn;" and -IR: escribir, " to write." When you conjugate a verb you take the generic, infinitive form and change the ending to match the person(s) or thing(s) doing whatever action the verb describes. In Spanish, the last two letters come off and another ending goes on in their place. Each verb tense has its own set of six endings for each verb type. Don't worry though, there are patterns there and much of the endings overlap so it is easier to remember them. The present tense verb endings are as follows: -ar verbs. That is: Sometimes the vowel in the ending is accented. Here is an example conjugation: Bailar Of course, most of the above statements can be restated without any loss of meaning with no pronoun: Bailo -- "I dance"; bailas, "you dance". In Spanish, the verbs change to reflect the person(s) doing the action on them and it is OK to leave off the pronouns, especially when context makes the subject clear. -er/-ir verbs. -ER and -IR verbs are almost identical, so they are often listed together. The only differences in the present tense are in Nosotros and Vosotros. For -ER verbs: For -IR verbs: That is: Here is an example conjugation of an -ER verb: Aprender Here's an example conjugation of an -IR verb: Escribir Irregular verbs. Tener. Note that the first "e" turns to "ie" in all of the conjugations except the "nosotros" form. That is a pattern that you will see repeated. 

Bicycles throughout the world are made with standardized, interchangeable parts. Unlike many modern products, the technology used in bicycles is simple enough to allow many riders to repair their own vehicles with a minimum of effort. For any cyclist, bicycle maintenance is a particularly useful skillset to acquire. Every skill learned in this area—no matter how simple or complex—can aid in keeping your bike in good working order, save you money, and make the difference between pushing your bike home or riding it. 

This page explains how to repair punctures in bicycle inner tubes, and gives some advice as to the most likely causes. In addition, there are notes peculiar to the use of tube sealants. Terminology. The basic parts of interest are these: The Causes of Flats. A "flat tire" or "puncture" is most often caused by glass, thorns, flints or nails when they cut through the outer rubber tread of the tire and damage the inner tube. These deflations are usually quick, at least when the object is removed. Although this is the most obvious way to get a flat tire there are other ways too. Consider these: Replace or Repair? It is undoubtedly easier to change a tube for a new one than it is to find the puncture in the original. For this reason many riders carry both puncture repair materials "and" a new tube when they travel for any great distance. When on the road it is sometimes difficult to find the point on an inner tube that is punctured. If it is a "slow" puncture it might be possible to pump some air into the tube, perhaps more than once, enough to reach home or a more convenient place to do the work. If is it a "quick" deflation, then a repair or replacement becomes necessary. When an object penetrates a tire it is best not to remove it immediately if it is the only thing that keeps air in the system. Instead, use the remaining air to get the bike home before removing it. If the tires are fitted with self-sealing fluids it might be that the puncture will have been fixed already and that the tire just needs to be inflated. A puncturing object should be removed from such a tire only when the part with the inclusion has been rotated to the "six-o'clock" position; this allows the internal sealant to fully reach the puncture. A different method is needed for slow-to-seal sidewall "pinches"; allow the fluid to pond at the bottom of the tire then tilt the wheel, slowly rotating it, so that the fluid can reach the internal sidewalls. If the tire is not self-sealing, and "must" be repaired rather than replaced, then the remainder of this page will explain how. The Basic Tools. This section lists the basic requirements for tube repair. The primary items are these: Also useful are these: The Repair Procedure. Examine the flat tire carefully to find any sharp item that may be responsible. If a nail, thorn, shard of glass, or a flint is found in the outer tire then mark the rubber of the outer tire with a wax crayon or suitable pen so that the location is easily found. Continue to check the tire in case there are more. Some punctures of inner-tubes also cause damage to the fabric of the outer tire, so look for any potential problems such as a bulge, that suggests a torn carcass, or for exposed carcass cord. If such damage is excessive then it may cause repeated punctures to inner tubes. At this time the method to correct such damage involves the surfacing of the inside of the tire with a so-called "boot"; a pre-glued patch that keeps the fibers clear of the tube. A piece of an old inner tube can also be used as a temporary fix, if loosely wrapped around the tube in the location of the damage. This latter fix suggests that it might be a good idea to carry an old piece of inner tube in your repair kit. At this time there is no rubber-based filler for cuts in "outer" tire cladding, so water ingress is likely. If there is no clear cause of the puncture and the location remains obscure, or if it is a "slow" puncture, then the wheel will need to be removed to get proper access to the inner tube. Methods are given below that include both wheel removal and repair with the wheel in place. Bike tire sticks on Rim - solution … It won't always happen but you can't get enough purchase on the edge to push the tire over the inner rim … So what do you do? After several attempts to lever off with the tire lever … The tire and tube were 100% empty of air … I put one foot on VERY FIRMLY on the tire (soft sole) and the applied 1 tire lever to push the bead over the rim and it came over the rim relatively easily. I was then able to access the tube. Inserting my spare tube, and making sure it was tucked snugly inside the tire (again completely empty of air), I again had difficulty coaxing the tire bead over the last couple of inches of the wheel rim. (sounds familiar ). by the way leaving the tire on the wheel rim I flipped it inside out and had a good look inside to try (and out ) to try and make sue there were no spikes, pins or thorns on either side that might puncture the new tube. I had applied all due pressure with both tire levers bring the tire levers closer together but the last few inches was eluding me. What did I do? You've guessed it. I removed one tire lever and once again I applied my trusty foot to on edge / rim of the tire, and then brought the remaining tire lever along the edge (both hands now available) until I got the tire bead to flip completely over the wheel rim. Eureka! Not to be recommended but any port in a storm! Remove the wheel? Wheel in place. Most punctures need the wheel removed, but if you are sure that you know "where" the hole is, you can do the repair with the wheel still on the bike. This method is popular on bikes that need wrenches to remove the wheels, and for rear wheels, even when they are of the quick-release type. However, if the "front" wheel is of the quick-release type, you will usually find it more comfortable to remove it anyway. The sequence for an on-the-bike repair is just: Remove the Inner Tube. Be careful too, with the valve when it is removed from its rim, and remember to first remove the locking screw on the rim if the valve is of the Presta type. Valve Problems. Any inner tube can be checked for leaks by first inflating it, and submerging it in water. Telltale air bubbles will seen emerging from the site of the leak. This process can be made more effective by forcing sections of the tire between the hands to increase the leakage. A slow leak can sometimes be caused by the valve itself. Sometimes a leak can be seen in a valve by wetting its various parts with a soapy solution and looking for bubbles. To try to fix a leaking Schrader valve, deflate the tire, unscrew the valve body with a keyed valve cap or valve tool, and examine the seat or rubber sealing ring for cuts or nicks, dust, lint, or fibers that prevent the valve from closing fully. Likewise check the valve seat and the bore of the valve stem. Clean if necessary. The valve body may be replaced. If no spares exist, an emergency fix can be had by inflating the tire with the valve sealed with silicon rubber, caulk, or cured epoxy resin. Obviously the tube must be discarded after such a process. If the Presta valves cannot be disassembled, then tubes with leaking valves must be replaced. Some Presta valves, for example, those made by Schwalbe and those of Bontrager have removable cores, but most do not. Such valves can be recognized since they have flat sections on them to allow slackening and tightening of the cores with a valve removal tool. Internal tube sealants can block both valve stems and cores. However, stems are easily cleared with a thin object after the cores are removed, and sealant can be removed from the cores by washing them in water. Self Sealing Tires. Sealants will seal most small punctures in tires up to about one eighth of an inch in size. The product ("Slime"), is already in some tubes bought from bike shops, but can be poured into existing tubes, provided that they have removable valve cores. There are a few points peculiar to the use of tires with internal sealants. 

Overview. A bicycle wheel consists of a central hub and a round rim, joined by a number of spokes. Spokes radiate from the hub to the rim, where they are anchored in a screw-thread attachment called a nipple. By tightening and loosening the nipples, it is possible to bring the wheel back into round (vertical true) or remove a side to side wobble (lateral true). Spokes may be arranged in a variety of patterns, of which three-cross, four-cross and radial are the most common. The pattern affects the strength, weight and characteristics of the wheel but is not particularly relevant to the process of truing. Tools Needed. To build or to maintain a spoked wheel, a spoke wrench is necessary. The flat-to-flat dimensions of a typical spoke nipple are too small to fit any commonly-available open-end wrench. Any good bicycle shop will, however, stock a selection of spoke wrenches, with price and quality being proportional. A spoke 'twiddler', or nipple driver, is a tool that may be useful in the course of assembling a new wheel from its component parts - hub, spokes, and rim. It isn't necessary - especially if you aren't building a wheel - but it can save some time. Similar to a screwdriver, but with an off-set blade that rotates freely in the handle, it is used on the slotted outer face of the nipple - which resembles a slotted screw. A spoke 'twiddler' allows for rapid assembly of the wheel - or installation of a nipple on a single spoke - and is designed to self-limit the extent to which a nipple can be threaded onto a spoke. A spoke wrench is used to bring tension to the spokes, and is applied to the square part of the nipple that protrudes inward at the rim. Unlike the 'twiddler', the spoke wrench is designed to turn the nipple when the spoke is under tension. This difference is what makes the spoke wrench essential. A truing stand is a purpose-built stand into which a wheel (rim, spokes, and hub, with axle) is installed during wheel-building, wheel repair, or wheel maintenance. A truing stand is very useful - possibly essential - for making hand-built wheels. Effectively, it is a rugged, precision-made jig for holding the axle of the wheel solidly in place. By extension, the rim thus has a steady reference in the axle in the hub held in the truing stand. The rim is, therefore, also found in relation to a caliper, or set of calipers, built into the truing stand, and against which checks are made for radial and lateral true (read 'perfection') during the course of building or adjusting a wheel. For extreme accuracy in measuring true, an optional dial indicator may be fitted to the truing stand, and a separate tensiometer may be kept to hand to test for proper tension on all the spokes (tension-balancing). Instructions for creating an inexpensive, but very accurate, truing stand are here. A dishing tool is used to measure the extent to which the axle juts out past the rim. Since a true wheel has the plane of the rim centered laterally between the points on the axle at which the axle is fixed to the frame, the offset of the rim from that anchor point on one side of the wheel should be identical to the corresponding offset on the other side. A dishing tool is used as a comparator: the offset on one side is measured using the tool, that setting is 'stored' in the tool, and the tool is applied to the other side of the wheel for purposes of comparison; the deviation of the second side's offset from that of the first's indicates the direction in which, and the extent to which, the rim needs to be moved to make the wheel true. However, a dishing tool is not strictly necessary if you have a good truing stand: if it's understood that a true wheel (abstracting from the particularities of the hub and the spokes) is symmetrical (i.e. the rim itself is, in some sense, 'centered' on the axle) then you can use the truing stand's caliper (or calipers). and an occasional flipping of the wheel in the stand, to serve the same function as a dishing tool. Tools, Punting. In most cases, especially when truing as maintenance on a wheel that has already been built, checking the tension on the spokes by ensuring that they all make the same tone when "pinged" with a fingernail will work fine. Checking the dish is also unnecessary on minor repairs, as most good truing stands will be able to give an accurate idea of how close the wheel is to centered, though it is probably a good idea to take the wheel out and re-settle it in the stand to double check. It's necessary to measure the distance if you are truing a wheel that is meant to have an offset, but this is very rare. In a situation where a stand is not available, the brake calipers on a bicycle can be used to measure true while the wheel is still attached, but this is less than ideal. Also, if no other option is available, a small adjustable wrench can be used instead of a spoke wrench or key, but extra care should be taken not to strip the nipples. A glass cutter made by Richards has one opening small enough that it can be enlarged slightly to make a passable 0 spoke wrench - or a 1 or 2, if that's what you need. If a rim has been 'pringled', or 'potato-chipped', into a saddle shape, by a lateral blow, it will be difficult to straighten by spoke tensioning alone - in fact, it is probably impossible to overcome severe rim warping this way. Often, however, it is possible to restore the wheel to a nearly-true state, even if it seems hopelessly warped. The procedure that follows should be applied as soon as possible after the trauma; leaving the wheel in it warped state for more than a few days will likely cause the wheel to cold-set - meaning it will retain the new shape in which it is left, and probably will have to be replaced. In order to effect the following quickie repair, you must have a properly inflated tube and tire still on the rim, and the rim itself should only be generally out of shape, with no damage to the rim wall itself. To wit: standing with your feet shoulder-width apart, with both hands, grip the wheel firmly by the rim and rubber at the point directly opposite the point on the circumference that is most obviously bent away from the principle plane of the wheel; bend slightly at the waist and position the point of the tire opposite on the ground in front of you, with the wheel at an angle of 45 degrees to the ground (possibly greater, depending on the denomination of the tire), and the bent part heading earthward; raise that point off the ground a few inches and let it fall again, ensuring that in the next step, only the tire - and no part of the rim - is going to strike the ground; without straightening at the waist, or bending over further, raise the wheel off the ground to about head or chest height, then forcibly bring the wheel down in such a way as to strike the ground sharply with the inflated tube and tire at the point opposite your grip; the rim should pop back to very nearly true. You can't undo damage done by striking the rim on the ground, so plan carefully, and rehearse mentally before committing to this repair. If the quickie fix above has improved matters, it's usually possible to improve the situation further with the usual techniques required of truing or maintaining a wheel. In the event that the above procedure doesn't bring the rim back to at least a rideable state, it may be necessary to replace the rim or the entire wheel. Note, though, that a rim can often be straightened, at least to some extent, by placing the rim (with appropriately loosened spokes) between two rigid objects (e.g. a pair of closely-spaced pipes, or between a door and its cross-wise bar-like handle) and, using a prying action, creeping up on the desired degree of 'flat'. Don't bend the rim more than necessary, work iteratively, check frequently on your progress, and be aware that you may need to overshoot your target slightly to allow for the tendency of metals to spring back from bending. It should be noted that repairing local rim damage with pliers or even a hammer and anvil is never advised - because less dramatic, more effective, methods are available which will not cause even superficial damage to the rim wall - the all-important braking surface in most cases. Instead of pliers neat, repair a flairing of the rim by sandwiching the rim between two thin, very flat pieces of metal (e.g. two cone wrenches) and placing this in the jaws of channel-lock pliers. An alternative method for repairing rim flairing is to use blocks of wood as both anvil - below the rim - and hammer - above the rim. If it seems necessary to use a hammer, don't strike the rim directly with the hammer's face; lay a strip of aluminum or brass over the area of the rim to be repaired, and strike there. Aluminum cans - ubiquitous and easily cut with scissors or even a pocket knife - are indispensable in bicycle repair. Truing the Wheel. First, remove the wheel from the bike and remove the tire and tube before placing it in the truing stand. Adjust the arm and caliper on the stand so that the caliper is just shy of touching both sides of the wheel. Now, spin the wheel and slowly close the stand's calipers until they scrape against a spot on the wheel. Once you've found a spot where you are out of lateral true, tighten the spoke or spokes on the side opposite the bump, and loosen the ones that are pulling it out of true. Be patient while doing this - you shouldn't be going more than a quarter turn at a time while truing like this. If the spokes are giving resistance, try over turning slightly, and turning back to where you intended. After you've gone through several passes like this, you should check the vertical true of the wheel to make sure that you haven't put it out of round. Re-adjust the arm and caliper of the stand so that the calipers are together and just underneath the rim. Spin the wheel, and this time adjust the arm until the wheel begins to scrape. While adjusting for vertical true, you should tighten the spoke at the center of the hump, and tighten the spokes to the sides one half as much each. Because these spokes will be on opposite sides of the wheel, this will ensure that you don't put the wheel very much out of lateral true. If the hump is between two spokes, tighten them equally. Adjusting the spokes in one place will affect another section of the wheel, somewhat like squeezing a balloon. After this is done, you should check for lateral true and even tension, retruing for both lateral and vertical true if the wheel is out. If the wheel is properly trued and tensioned, you should stress test the wheel by placing it on one side and pressing down on it fairly firmly. You should repeat this going around the wheel, in order to be sure that the spokes settle into position with the spokes that they cross now and not while the wheel is being ridden. This can also put the wheel out of true again, and this should be checked. A full, comprehensive discussion of bicycle wheel building and truing is found in Jobst Brandt's book "The Bicycle Wheel" . 

 ADJUSTING RIM BRAKES Rim brakes on bicycles are simple to adjust, and the term "rim" is used to distinguish them from "hub brakes". Rim brakes include all brake designs that depend on using brake pads to close on the rims of a bicycle's wheels. The brake parts for the most common configuration are identified in Figures 1 and 2. Although brakes are usually mandatory on bicycles, the laws and rules for brake performance vary, and some countries, while making law for "new" bicycles, ('at the point of sale'), have few specifications for the operation of bicycles "after" that point. Bike riders are advised to use the sources of law that apply in their "own" countries, but where the rules there are incomplete or unclear, they are advised to use any "point of sale" specifications that are available. Because the matter might then still be unclear in some countries, a table below repeats the braking distances of British Standard BS6102/1 for new bicycles. Bike shops can perform any number of tasks for the bicycle owner, but the basic brake adjustments are easy to do for yourself. At the simplest level they consist of screw adjustments on the handlebars, while knowing what gaps you intend to produce. At other times, (rarely), the cable length needs to be changed, but this too can be done by anybody with a practical leaning. It is perhaps when the conventional adjustments fail to solve the problem that most people resort to such pages to learn more. This page includes a selection of the most common confusions for brake adjustment and explains how to correct them. Bear in mind that the best way to learn the adjustment of brakes is to be shown by somebody while it is being done; in this way much of the mystery vanishes. The next best way is to follow a fairly stolid description of the sort below, and although it lacks interaction, should at least leave the reader better informed than when he started. &lt;br clear=left&gt; Preliminaries.  Consider these few things before carrying out brake adjustments. Doing so might save time in the long run: The Full Procedure The full adjustment procedure is summarized below but it should be emphasized that slight adjustments might solve the problem. In any case, it is quite usual to repeat balancing at various stages throughout the adjustment. These above points are all described in some detail in the text that follows. Brake Block Alignment. The brake blocks need to be aligned with the metal rims. Refer to Figs 1 and 3. The leading edge of each block should be slightly closer to the rim than its trailing edge. This prevents brake squealing and is called "toeing-in". Use a coin, a credit card, or any other thin material under the back end of the block while adjusting it. Some suggest tying an elastic band temporarily to the trailing end of the block to allow more freedom while working. To make the adjustment, slacken the screw that holds the block. Usually a 5mm hex wrench is used. Swing the brake arm in so that the block is pressed squarely against the metal rim, and then re-tighten it while holding the block hard with your "toeing-in" device in place. Avoid the rubber of the wheel; the block should contact only the rim. The block should be parallel to the rim, noting that some blocks are curved to fit its shape. Do one block at a time, and just let each arm relax after the block is set. If there is insufficient clearance to work or if you intend adjusting the cable length later in any case, then unhook the cable bridge (Figure 1, D) or undo the cable clamp (Figure 1, F) before carrying out the work. Block Clearance. Decide whether or not the block clearance is correct by trying the feel of the brake lever. (Figure 2, D). The brake should feel responsive without too much handbrake slack prior to the start of braking. Some mountain bike V-brakes might need only a 1mm gap, while many other brakes need about 2mm. If in doubt, refer to your bicycle handbook. If a "significant" adjustment is needed, resetting the cable length should do it. If a "small" change will do then use the brake lever barrel-adjusters on the handlebars, as described in the Fine Adjustments below. Rough Adjustment. Alternatively, to set cable tension, back the barrel adjusters off a few extra turns, squeeze both pads against the rim, tighten the cable pinch bolt, then back of the barrel adjuster until the wheel spins freely between the pads. Brake Block Balance. 'The brake arms should be adjusted so that both blocks apply pressure to the rim at the same time. As a result, at balance, there is no sideways displacement of the wheel during braking. Although slight imbalance is not "always" critical, displacement of the wheel by even a small amount can cause damage when small-clearance devices such as distance counters are installed on the spokes. Balancing the spring tensions keeps the wheel centered even during braking. For brakes like V-brakes, there is a small screw near the bottom of each brake arm to adjust the spring tension.( Fig 5). It is often a posidrive screw, (M4x6mm) , with a tightening insert. Turning this screw "clockwise" will cause the brake block to move "outward" slightly, and turning it "counterclockwise" will cause the block to move "inward". As one block moves, so does the other, to maintain the distance between them. Adjust these until the clearances are about equal. In this way, operating the brake causes the blocks to reach the rim at about the same time. Be careful not to withdraw the screws too far since they may not be captive. At the other extreme, if a screw is too far in, the brake arm will bind; if a brake arm seems inactive, or unresponsive, this might be the case, or the spring may just have popped out of its slot. Try to avoid the limits and to reach a balance with the screws near their mid-range. This is easier than it sounds since making an identical adjustment on both screws will leave the balance unchanged. Brakes usually can be balanced unless the wheel is not centered in the wheel-arch. (See Common Brake Problems on this page for more on this). Final Check. Rotate the wheel to check brake clearance. Make sure that there are no repetitive noises coming from the brakes. Test the brakes on the spinning wheels before riding. If these work well enough then test the brakes again by riding the bicycle in a quiet place. The brakes should stop the bicycle decisively in a fairly short distance. Some v-brakes in particular have a short stopping distance; on these you should not need a deep pull on the brake lever for a good braking effect since this is a sign that the blocks are set too far from the rims. In any case be sure to refer to the manual if there is doubt. Figure 4 is an extract of maximum braking distances as given by the CTC Hire Standard, that is itself related to the content of British Standard BS6102/1 for new bikes. The CTC standard is an attempt to consider used bikes, as opposed to bikes at the point of first sale. In any case, these stopping distances are useful until such time as the European standards properly address the issue. The most common reason for long braking distances, apart from maladjusted brakes, is the degradation of the brake block surface area. Be sure if replacing these to replace both together. Common Brake Problems. Balance Problems. Sometimes, despite best efforts, the brake arms will not balance. The spring balance screws are designed to have limited range since they are only expected to handle the "difference" between the arm tensions to achieve balance. So, faults that only slightly bind any part of the braking system can cause trouble with balancing. Possible faults include: The lack of general lubrication, the binding of brake arms, the slipping of faulty adjustment screws, unseated brake-arm springs, unseated housing ferrules or nipples, or a quick-release wheel that needs re-clamped closer to the centre of the wheel arch. See these points below. Soft Brakes. If brakes are softer than intended, even after adjustment, then it might be that the toeing-in is excessive. Also, if wheel wobble causes rubbing on the brake blocks, widening the gaps will necessarily soften the brakes. See the comments below. Cables and Housings. The subject has been given a separate page at Cables and Housings. 

Mending a broken chain. Newer, narrower chains and wider chains for single speed or three speed bicycles often have special links for removing the chain, but the chain tool is still needed to remove excess links when replacing chains. A chain tool can be used to push out the rivet which joins the plates of the chain together. Many chain tools have two positions the chain will fit- make sure the plate furthest from the extractor pin is supported by the tool. Once this is done the chain can be replaced or re-assembled as a shorter chain. Most chains will have enough slack to allow removal of a few links without making the chain too short. Avoid pushing the rivet all the way out, it should remain lodged in one of the outer plates. A small portion of the rivet left protruding on the inside of the plate can hold the chain together during reinstallation. Use the same position in the chain tool, with the chain plate furthest from the chain tool's pin supported, and push the rivet back into place. After reinstallation, the link will usually be too stiff- leading to chain skip unless loosened. The other position in two-position chain tools, with the closer plate supported by the tool, is used for loosening tight links. Force the side of the rivet that protrudes more into the link to loosen the link. Alternatively, grasping the chain on both sides of the tight link and flexing the chain in and out will loosen the link. Note. Although they're made of metal, chains seem to "stretch" as the links wear. This will cause them to not engage on the gears properly (they will hook on to the middle or top of the tooth, instead of the bottom of the groove). This will also cause the gear teeth to wear down. To test for excessive wear on a chain, open a link and remove the chain. Then try to flex the chain sideways, in the direction it is not supposed to bend. If you can make anything more than a 1/8 arc (for example, if you can make a half circle or if you can touch the ends), then your chain is worn and should be replaced. Alternatively, measure a straight section of chain under slight tension. Standard links of chain measure one inch when new- if 11 links of chain measure 11 1/8" (283 mm) or more the chain should be replaced. If you only replace the chain and not the gears, the chain may skip. For this reason, it is best to replace both at the same time (the front gears are not as affected by this because they are bigger and the chain doesn't pass over each tooth as many times). Replacing the chain more often, when 11 links measure 11 1/16" (281 mm), will allow the gears to be reused . Avoiding Chain Wear. To prevent such wear, avoid using gear combinations that stretch the chain diagonally. Also change gears only when you are turning the pedals. Remember to gear down when coming to a stop. 

Regular bicycle maintenance. Performing regular maintenance on a bicycle will improve its performance and longevity, and reduce the risk of breakdowns. The exact schedule for a particular bicycle will depend on how it is used: its weekly mileage, the weather conditions, road (or off road) surface conditions and so on. Most parts will need attention and possible replacement every year or two; if this is done, however, a bicycle can be maintained in good working order for decades. Bicycle inspection &amp; maintenance can be roughly broken down into four categories. Each includes clean, inspect, adjust, lubricate and repair as necessary. The primary difference between them is the depth or level of each task and sub-task. The schedule given here is a starting point for an average bike, assuming daily or weekly use; you will soon adjust this based on your own experience. Every ride: Once a week: Quarterly: At least every two years: 

Begin the Wikibooks Dutch Language Course! 

=Welcome!= Welcome to the Dutch language course. Notice the arrows under the images of a number of cities where Dutch is spoken? Click on the arrow to hear how their names are pronounced in Dutch. There may be some surprises! Information about the course. If you want more information on the course, here are a few pages that may give you that. The lessons. If you want to jump right in and start learning: there are three types of lessons: There are three levels, each consisting of two cycles of four lessons each. Beginner level. At the end of this level learners should be able to form simple sentences in the basic tenses and possess a vocabulary of just over 1000 terms. Intermediate level. At the end of this level learners should be able to deal with complex sentences with complex verbal expressions and have a reasonable grasp of syntax. Advanced level. At the end of this level learners should have full command over Dutch grammar and syntax, including a number of special topics. Under construction. /For children/ 

Layout of the Course. This textbook is intended to be a comprehensive course in the Dutch language for English speakers, but of course people who speak English as a second language are most welcome as well. Being an Afrikaans speaker is a huge help too, as it is a daughter language of Dutch, though Afrikaans has its differences and words can have different meanings from the same ones in Dutch, and about 95% of Afrikaans vocabulary comes from Dutch. If you are an Afrikaans speaker, this course shall be significantly easier. Just remember the different dialects etc. when speaking Dutch. Early lessons emphasize conversational subjects and gradually introduce Dutch grammatical concepts and rules. In addition, sound files and illustrations accompany appropriate parts of each lesson. Structure of the course. The course is divided in three levels: beginner, intermediate and advanced. Each level consists of two cycles of 4 lessons each. At the beginner and intermediate level each lesson is accompanied by two parallel lessons: a "practice" lesson and a "cultural" lesson. That means that the beginner level comprises 2*4*3= 24 lessons in total. At the advanced level there are no parallel lessons. Starting from scratch, expect to arrive there in a year, depending on how intensively you study. By that time you should be able to start picking your own things to read, newspapers, Wikipedia articles, what have you. Parallel lessons. The main lessons aim at introducing grammatical topics by means of conversations, interspersed with some exercises. Of course that is not sufficient to actually start speaking the language. Therefore, each lesson is accompanied by: In addition there are pages intended to help build up vocabulary, some of which are interwoven in the practice lessons. Other ones are stand alone. Possible strategies. Which way the reader wishes to use the book may vary. The recommended strategy for a beginner with no experience is: Lesson 1 &gt; Lesson 1A -&gt; Vb. 1 -&gt; Lesson 2 &gt; Lesson 2A &gt; Vb. 2 &gt;Lesson 3 &gt; Lesson 3A &gt;, etc. People who have experience with other languages, grammars etc. might want to follow the order Lesson 1 &gt; 2 &gt; 3 &gt; 4 &gt; and on to the end of the basic text Others may want to start tackling the language in context of a situation and worry about grammar later might want to start with Lesson 1A or Example 1 and check back for the grammar in Lesson 1 later. But readers are encouraged to revisit pages they have worked on again and again, so the order may be more complicated than this. For people who have been learning Dutch in other ways and want to revisit a certain topic there is an Index of grammatical and syntactic topics covered. Layout within main Lessons. The following topics can typically be found in one of the main lessons: Layout in the practice lessons. There are often additional texts, many exercises of various sorts and short quizzes to practice what was taught in the main lesson or to repeat what was taught in the lessons before. To further expand vocabulary students may be asked to go to one of the vocabulary pages to study words related to a certain topic. Quizlet links provide another way to brush up on vocabulary Layout in the example lessons. In these lessons we will have a look at rhymes, poems, songs etc. Often the student is referred to a YouTube video to watch in order to practice oral understanding and expand vocabulary in the context of Dutch culture. Quizlet links provide another way to brush up on vocabulary. Quizzes. Short quizzes are integrated into the practice lessons, but there are also longer quizzes at the end of the two cycles of the Beginner Level. Rate of acquisition of vocabulary. Each lesson -its two parallel lessons included- contains on average about 120 new terms. At the end of the Beginner Level the student should know a little over one thousand terms. The Student and the Lesson. The text is designed to constitute a comprehensive course of study in the Dutch language. Each lesson should be read thoroughly and mastered before moving on. Substantial text in Dutch is included and the student should read all of it, not once, but multiple times. In most cases sound files are given as an arrow button; they should be listened to multiple times as well, both while reading the text simultaneously, and -once the content is understood- with eyes closed. Complete translations into English are included only in selected places or they are hidden in a drop down box. Most of the text should first be translated by the student using his or her acquired vocabulary and the vocabulary presented at the bottom of each lesson and/or in collapsed boxes in the right hand margin of the text. Hints about the meaning of new, underlined words can also be obtained by the . As the Dutch is read (out loud is better), the student must succeed in gaining an understanding of the meaning of each sentence, and the role each word plays in establishing that meaning. To the beginner, there will seem to be many words in a Dutch sentence that are out of place or even redundant or unnecessary. These add subtleties to the language that will make sense eventually. But it is important to experience these subtleties from the very beginning. There are exercises interspersed throughout the main lesson with additional ones in the practice lessons. They are of various nature. Translation exercises, pronunciation drills, fill-in-the-blanks, adapt a word form, swap to items etc. The Dutch Language. Dutch ("Nederlands") is a member of the western group of the Germanic languages. It is spoken primarily in the Netherlands, and in a major part of both Belgium and Surinam. It has about 23 million mother tongue speakers and another 5 million second language speakers. Continue reading about the Dutch language and its history at Wikipedia. Dialects. As a standard language Dutch is a relatively young phenomenon. The standard is based on a variety of dialects that are much older and show considerable differences not only in pronunciation but even in grammar and syntax. This holds for many languages, including for English as spoken in the UK. By urbanization, suburbanization and the influence of the mass media the standard language has been gaining ground at the cost of the dialects for over a century, so that it is now the mother tongue of most speakers. Others are typically perfectly bilingual in their regional tongue and the standard language. But in the way that the standard is spoken there are many regional and social differences in pronunciation or even in syntax and grammar. In Bruges (Flanders), Rotterdam (Netherlands) or Paramaribo (Suriname) Dutch will sound as different as English does in Edinburgh, London or Indianapolis. This course aims at teaching Dutch that would be acceptable to most if not all speakers but will point out a number of important differences that non-native speakers are likely to encounter in their interaction with native speakers. A major division in the dialects is formed by "de grote rivieren" the great rivers that run through the Netherlands from east to west on their way to the North Sea as shown in the map. The course is mostly based on Northern standard Dutch, because that is most readily accepted under all speakers, but it will point out some important differences at times and also give examples of other varieties. A dynamic language. Dutch has undergone far more sweeping changes in grammar and syntax in the last century or two than either English or German. It has lost most of its case endings and much of one of the three original genders (feminine). This has led to some interesting shifts in its grammar and syntax. Some of these developments are still taking place today. This means that Dutch grammar is less set in stone than the reader may be familiar with from other grammars. Occasionally we will have to discuss the evolution rather than the creature to explain modern Dutch usage. Dutch and English. If you are an English speaker unfamiliar with Dutch, you may be surprised to learn that English and Dutch are closely related languages and share many words that are very similar. This is particularly true for everyday words in English that are Anglo-Saxon (i.e. Germanic) in origin. After 1066 English has absorbed a lot of (Norman) French. Dutch also has been exposed to contact with first vulgar Latin and then French, but the French influence has been less pervasive. Consider the following list of English words followed by their Dutch counterparts: Many words of French origin have entered both languages and are quite recognizable: But in many cases Dutch retains a Germanic word, sometimes aside the Latin one: English spelling has conserved many now silent consonants, e.g. "gh" in "light". This may have been an obstacle when learning to write English but when learning Dutch the investment pays off. Dutch has "licht" and the "ch" is very much still pronounced as a guttural fricative /x/ like in German Bach or Scottish Loch. Some words are even completely the same. ("true friends") Of course, even words whose spelling is no different in English and Dutch may be pronounced quite differently or mean something different ("false friends"): Nevertheless, when reading Dutch you will see the kinship between the languages, even in many short words, common or not. For example, compare: These sentences consist almost entirely of "cognates": words that evolved from the same source. Dutch is indeed one of easiest languages to learn for an Anglophone. Notice however the position of the verb is in these two phrases. In Dutch it stands in front of the father. This is because Dutch has retained something that English has lost: the rather complicated word order (syntax) of the West-Germanic languages. Many English speakers who learn Dutch find that one of the most difficult aspects to learn to do correctly, but it hardly ever leads to miscommunication. In the course it is introduced bit by bit. The full picture is only described in one of the last lessons (21). Dutch and German. Both Dutch and German are West-Germanic languages and this means that there are many resemblances. However, Dutch is easier to learn for a speaker of English for a number of reasons. First of all, (High-) German underwent a major shift of almost all its consonants in the early Middle Ages. In term of its consonants Dutch has been pretty conservative. Compare: This makes a major part of Dutch vocabulary easier to memorize. Secondly, German retained its system of case endings in contrast to Dutch and English. It is not easy to master that system if your mother-tongue does not have it. Compare: Knowledge of German can certainly help in learning Dutch, but it can also be a source of confusion. A good example is the letter combination sch. In German it denotes the same consonant as sh in English (in IPA: [ʃ]), in Dutch this sound is relatively rare. It only occurs in loans from languages like Frisian, English, French etc. In Dutch 'sch' usually denotes [sx]: an [s] followed by a velar spirant [x], like in "schip". In the ending -isch the 'ch' is mute and it is pronounced as [-is] as in English 'fleece'. A topic where knowledge of German is a great help is syntax (word order), but on the other hand there are differences in how the verb tenses are used. In German the imperfect past tense, like "du gabst, sprachst, rettettest" is on its way out. In Dutch forms like "jij gaf, sprak, redde" are alive and kicking and in every day use. Dutch and Scandinavian languages. Although Danish, Swedish and Norwegian are North-Germanic languages, which means that the relationship is a bit more distant, speakers of these languages typically do not find Dutch very difficult, because many changes in the language such as the loss of the case endings and of feminine gender are quite similar. Some Scandinavians manage to speak Dutch so well without any discernible accent that people don't realize that they are not Dutch. Swedes may have to pay attention to their intonation, because that is rather different in Swedish. Dutch and French / Italian and other romance languages. For speakers of the Romance languages Dutch is by no means an easy language to learn, although if you already speak English some of the problems may already have been overcome. Dutch is a stress language like English and German. That is: every word has one syllable that is high in pitch, a bit louder, and usually a bit longer than all the other ones. In French those three qualities are not coupled and spread more evenly over the syllables of the word. In Dutch, stressed syllables have either a full vowel (one of twelve) or a diphthong; the unstressed syllables all tend to have a schwa. Intonation is therefore difficult for Romance speakers because the contrasts between the syllables is much smaller in those languages. Another problem is formed by the separable verbs. There often is no direct equivalent in French for the fine nuances imparted by the separable prefix in Dutch. Often an entirely different verb needs to be substituted, or there is none available. French speakers usually have little problem with the vowels, but they do tend to speak much more in the front of the mouth than Dutch speakers do. The placement of the sounds is different. For other Romance languages a vowel like "u" [y] or "eu" [ø] the "ui" [ʌy] diphthong might be problematic, as well as the guttural spirants "g" and "ch". Dutch and Russian / other Slavic languages. A main problem for Russian speakers lies in the vowel system of Dutch. The vowels a, e, i, o, u all occur in two different qualities. The differences are difficult to hear and reproduce for Russian speakers. Another problem is the voiced consonants like b, v, d, z. In Dutch they all get devoiced at the end of the word, but that hold for Russian too. However, they often also get devoiced when they follow a voiceless consonant, even in assimilation from the previous word. E.g. A grammatical problem is that Dutch is a "tense" language and Russian an "aspect" language. Dutch has perfect tenses, Russian perfective aspects. They cover the same territory, but one lengthwise the other in the other direction. Dutch and Mandarin. Of course the differences are very large and numerous, but one difficulty deserves mention, that of epenthesis. The syllables of Mandarin are either CV or CVn, i.e. they start with a single consonant and end in a vowel or a nasal. In Dutch they can end in other consonants and in fact multiple ones, like "herfst" is CVCCCC with no less than four consonants at the end. Mandarin speakers need to suppress the tendency to add epenthetic vowels, making it something like herefesete. Epenthetic (inserted) vowels are sometimes used by Dutch speakers too, e.g. "werken" may well be pronounced "werreke", but this is considered dialectal and non-standard and is frowned upon by most speakers. Vocabulary and Grammar. In learning to read or speak any new language, two important aspects to be mastered are "vocabulary" and "grammar" (others are pronunciation and syntax, but they usually do not stop you from being understood). Acquiring vocabulary is a "simple" matter of memorization. Learning by ear. Children do it all the time, but they are at an advantage: they memorize far easier than grown-ups. Age is a definite disadvantage in language learning. The child's learning process can be "reactivated" to some extent by immersion in a second language: a method of learning a new language by moving to a place where that language is spoken and having to get around and live without use of one's native tongue. If you do not have the opportunity of residing in a Dutch speaking area an alternative is to listen to recordings and we are in process of adding bits and pieces as .ogg files so that you can learn by ear. Use them as much as you can. More than once. These files take different forms In the Example lessons there are also links to YouTube videos where the text of a poem is recited or a song is sung. Of course there is also a drawback to the by-ear method: You do not get much immersion into reading Dutch. You as an internet user, will most likely want to be literate in Dutch. As with all languages: the "written" Dutch language and the "spoken" Dutch language are by no means identical. At times we will explain the distinction if necessary. Learning by eye. This is why this course also tries to train your eyes, but this will not work without effort from your side. This is why we often say: Your turn! (Uw beurt!) So what do you need to do? There are a variety of things. We are tackling the problem with a multi-pronged approach. Be sure to "learn"—commit to memory—all of the vocabulary words in each lesson as they are presented. Early lessons have simple sentences because it is assumed that the student's vocabulary is limited. To help you accumulate vocabulary there are a number of additional pages see: Dutch/Vocabulary. They are mostly both visual and audio in nature and there are exercises to go with them (still being created). Throughout the text, more complex discourses (e.g. as photo captions) are included to introduce the student to regular Dutch in use. It may be helpful to translate these using a Dutch-English dictionary (Wiktionary is usually on a click or two away). Other sources of Dutch, such as newspapers, magazines, web sites, etc. can also be useful in building vocabulary and developing a sense of how Dutch words are put together. The Dutch Wikipedia provides an ever-expanding source of Dutch language articles that can be used for this purpose. Further, a Dutch version of the English Wikibooks project—a library of textbooks in Dutch — is available at and there is a growing Dutch version of (WikiWoordenboek) to which a number of words in the text have been linked for direct reference. WikiWoordenboek usually has an example phrase to go with every dictionary entry to show the word in context. This too is a helpful tool for expanding your vocabulary, as context helps memorization. Learning grammar and syntax. This is where as a grown up you are at an advantage, because you may already know how grammar works from your mother tongue or other languages you are proficient in to some extent. Dutch grammar is sufficiently similar to English grammar that "reading" Dutch is possible with minimal vocabulary. The student should generally recognize the parts of a sentence. With a good dictionary, a sentence can usually be translated correctly. Of course there are some notable exceptions and "false friends", e.g. in the way that the passive voice is formed: To speak and write Dutch properly you do need to learn its grammar and syntax. Particularly the latter (word order) is rather different. We will gradually introduce it. Do not be daunted by it. Learning a language goes bit by bit, word for word, structure by structure. Just keep at it and look at what you have gained not at what you don't understand. Children don't always understand everything either, but they are not ashamed or humiliated by that. Pronunciation. A guide to pronunciation of Dutch is provided as Appendix 1. You should become familiar with this page early on, and refer to it often. Nothing can replace learning a language from a native speaker, but the text is liberally sprinkled with audio files providing the student with valuable input from hearing spoken Dutch. Analyze the spoken words carefully. Descriptions of the pronunciation as in Appendix 1 can only closely, not exactly, convey how Dutch words should be pronounced. This is why there are quite a few buttons on that page that allow you "hear" what is meant. Do compare! Of course there are variations in pronunciation. Dutch spoken in the country side of Brabant, in the cities of Amsterdam, Antwerpen or Paramaribo sound pretty different and the same thing can be said for the board room and the back alley. We do present pronunciations from different people and places in the sound files and even more so in the YouTube video links. Contact. If you have questions or want contact try the dutchgrammar forum. I often hang out there and correct people's stuff or do some exercises 

Grammatica 1-1 ~ Grammar versus what children do. Why grammar? Children learn their mother tongue without knowing the parts of speech such as verbs, nouns and phrases. However these are helpful for anyone attempting to learn a second language from a book or a website. Of course the children have it right: the "best" way to learn a language is to listen to a mother tongue speaker and simply repeat. Then just use the word in a similar situation and see how people react. Children are masters at acquiring language this way and are generally smiled at when they use a word incorrectly. Being an adult, people are often not so forgiving to you and you feel foolish when people laugh and point out to you that you just said "toothbrush" while you meant "toothpick". Besides, native speakers may not always be available to you. Or if they are they are not eager to spend time playing 'child' with you. This book will try to compensate this by addition of audio files and visual information, -as the figure to the right- but that is still only a cumbersome substitute. We do recommend that you use them as much as you can. Firefox seems to give easier access to them than other browsers. So, please go ahead, push that arrow and learn to say 'toothbrush' properly. Your first Dutch word? Congrats. Although clearly children are superior in language acquisition, grown ups do have an advantage: they can analyze language better in terms of its grammar. This course therefore uses both approaches: it will deal with grammar, but it will also ask you to be a child and listen and repeat or look at some pictures while playing a sound clip. Don't be afraid to be a bit childish! It serves a purpose. One important observation about children should be mentioned: they always learn language "in a certain context". Is it all just grammar here then? No! There is much more. Audio files are inserted into the main lessons as much as possible, even though they aim at gradually introducing grammar and syntax. The parallel series of practice lessons (1A, 2A etc.) provide additional practice, vocabulary building, sound material and quizzes. The example pages (Vb. 1 etc.) follow the contextual path of learning like children do and involve nursery rhymes, poems, stories, songs and the like. And there are audio-visual vocabulary pages to help you learn more words. What is the best way? So, what is the best way to learn a language? The best way is to do something everyday. What you do is often less important than simply doing it. Children are champions in language acquisition and they never worry about what they do. Oh, and what you do, may very well be doing that same exercise again. Children love doing things "again". Ever watched the Teletubbies? "Repeating" is an important key to language acquisition. Being "efficient" and saying: "oh, I have done that before, let me skip that!" is a bad adult habit that children would never stoop to, until they get really bored with something. (Which is when they already know it). So, push that button below the toothbrush again! And tomorrow come back here and do the same. Other assets. Another thing to exploit is the other languages you already know. English speakers will find many strong parallels between their language and Dutch. German speakers even more so, but there are also differences. Where possible we will try to point out the similarities and the differences and exploit them. However, as noted in the introduction, Dutch grammar is more complex than English grammar, and identifying the meaning of words in a Dutch sentence is difficult without understanding the clues to word function that come from the grammatical rules. The basic lessons of this textbook are set up to first introduce the parts of speech, and then bring in the rules that govern these. Pay particular attention to sentence word order as you progress through the lessons. Some tricks this course uses. Hovering. Some words will be . Try to hover your mouse over such words. Topics and vocabulary. There are pages to help you build vocabulary in a visual / auditive way. Audio files. Whenever you see one of the following: Or: Please click and listen! (If you do not see any buttons now: try a different browser. Firefox and Chrome seem to work. Internet Explorer does not.) After listening, pronounce the word the best you can and then click again. Keep doing that till you are satisfied with your own result. It is useful to then leave it be for, say 20 minutes and do it again. Then perhaps once more the next day. Vocabulary / Pronunciation boxes. For many texts there is a box on the right that you can open to look at the vocabulary being trained and listen to the pronunciations. Click it to open it and start listening and reading. Using other sites. Wiktionary. Throughout the texts and in the vocabulary lists there are blue links that take you to the Dutch version of our sister project Wiktionary. It is called WikiWoordenboek. Of course the layout is in Dutch and you may not immediately understand everything, but that is not a disaster. If you want to learn a language you also should learn to be a bit of a detective: you often need to get the gist of something with a few pieces of the puzzle missing. Don't let that scare you off! Here are a few useful topics used on WikiWoordenbook: If you are really lost use the interwiki link to the English version (or any other language you know) as back up, but don't give in to it too easily! Use it to figure out what you did not quite get on the Dutch version. We encourage you to use the links to expand your vocabulary. First guess what a word means, then click! Just try it on the verb ; look at the box with three forms on the right. What is the past tense? (It is the middle one in the box) Quizlet and Memrise. There are plenty other sites that allow you to expand your growing knowledge of Dutch. They all have their pros and cons. For example Memrise and Quizlet have an interesting way to boost vocabulary, but teach zero syntax or grammar and usually little other context. But if you want a vocabulary boost it's great and we are in the process of creating practice sets dedicated to the material of the lessons here. Some already have a quizlet link on the bottom of the page. It is therefore recommended to register for Quizlet. (It is for free). YouTube. YouTube has a plethora of videos, that we will even send you to at times. Again: good for listening, vocabulary building, often less so on grammar and syntax, but plenty of other context. But do come back: the context. There is a pertinent Dutch proverb: "Verandering van spijs doet eten" - "Change of food makes you eat". Children also vary what they play with. Variety of learning is healthy. Besides children never learn a word outside context. When they learn a new word, it always goes together with the context of: One aim of this book is to provide as much of that context as possible. This is why conversations, stories, texts, poems and songs are important: they give context to the words you are learning. So, enough talk! Let's get started. And we'll start in the context of a simple conversation Gesprek 1-1 ~ Vrienden: Jan en Karel. We will put such text material in a colored box. What you are supposed to do with such material is the following. When learning a new language it is very important to be able to deduce meaning from limited information, because you will often not know all the words used. Picking up their meaning from context is an important skill. This is why the hovering is important. You may notice that Dutch sometimes strings words together a bit differently than English. Dutch word order is quite different and a difficult aspect of the language, but we will revisit that many times. So don't worry about it for the moment, just observe. Pronunciation. Dutch pronunciation varies with region and speaker, and you may have been shocked at some of the sounds of the language. You can visit Dutch/Alfabet if that is the case. Dutch spelling is not really phonetic, but pretty systematic (much more so than English) and once you learn the system you should be able to pronounce an unknown word on sight pretty well. It is not easy to render the sounds in writing, but the following rendition in IPA gives a reasonable idea. Try running the sound file again while reading the IPA version. If you prefer other renditions than IPA try this page Grammatica 1-2 ~ Forms. We will use the material in the colored boxes to point at grammatical phenomena and introduce the grammar that way, step by step. Clitic forms. Did you notice the difference between "Hoe gaat het met "je"? and "En met "jou?" in the conversation? Both translate literally into with "you", but there is a difference in emphasis. "Jou" carries emphasis, "je" does not. In Dutch, there are often two forms of the same pronoun: a strong one and a weak ('clitic') one. This is particularly true in spoken, colloquial Dutch. In the written language the clitic ones are not always shown as such. In colloquial English the same thing can be heard at times: "seeya!" instead of "see you!". In Dutch the use of clitics is very common; it already was in the Middle Dutch period before 1500. For now remember: never stress a clitic Polite forms. The above conversation was between two good friends. It utilizes the familiar form of the personal pronoun ("je", "jou") where English uses "you". However, Dutch also has a polite or formal form of the personal pronoun for the second person (you), u. Many languages have this distinction. It is e.g. comparable with in Gàidhlig, Sie in German, vous in French, usted in Spanish, Вы in Russian, or anata in Japanese. When to use one or the other is not always easy to decide. Someone unknown, particularly if older, is generally "u", an old friend typically "je, jou". The latter roughly corresponds with the 'first name basis' in English. Notice the use of "u" in the conversation below that takes place between colleagues rather than close friends. They would never say: "hoi!" to each other. Regional forms. In the South of the area where Dutch is spoken (Flanders mostly), , the familiar form when speaking in familiar fashion is "gij" (clitic: "ge", object: "u"). "Gij" is mostly used when speaking dialect, although it gets used more annd more in polite situations and on tv. In the north it has become obsolete since about 1800. It is used much like "you" in English for both singular and plural. In the North "gij" is now only encountered there in archaic phrases like: "gij zult niet stelen" - "thou shalt not steal". Like "thou" the pronoun "gij" takes its own verb forms. This course is mostly based on northern usage as this is the most widely accepted, including in Suriname and the Antilles, but some important differences will be pointed out and we will see "gij" occasionally when we look at some older poetry. Gesprek 1-2 ~ Collega's: De handelaars. Push the button and listen to the following text. It is recommended to first just listen. Then read the following conversation. It is a bit more formal than the one before. If you are not sure of the meaning of a word, hover your mouse over it, if it is underlined. A translation will pop up. Or use the vocabulary box to the right. Make sure you know the gist of the story before opening the translation box. If you cheat, you cheat yourself... Finally, listen to the recording again with your eyes closed. Can you understand what is being said? &lt;br clear="all"&gt; Go back to the pronunciation, close your eyes and see how much you understand now. You may have to repeat the process a few times. Quiz. How are you doing so far? Do this little quiz to find out! &lt;quiz display="simple"&gt; -toothpick -friend +toothbrush -visit -me -sir +you (object) -you (subject) ---+- already +---- street ----+ merchant -+--- good --+-- how &lt;/quiz&gt; Grammatica 1-3 ~ Introduction to pronouns. A is a short word that takes the place of a noun previously mentioned in the sentence, paragraph, or conversation. Recall: Kent u "meneer Standish"? Bent u "hem" al tegengekomen? "Hem" refers back to "meneer Standish". It is a "pronoun" that stands "for" (pro- !) meneer Standish. There is a variety of pronouns like personal, possessive, relative and indefinite ones. Let's look at the personal pronouns first. Personal pronouns. Both English and Dutch have had a system of case endings in the past, as languages like German and Russian still do today. In English most of the system fell into disuse starting with the Viking invasions of the 9th and 10th Centuries, and especially after the Norman invasion in 1066. The collapse of the system in spoken Dutch dates mostly from the 16th century and in the written language it was scrapped as recently as 1947. That means that Dutch has more remnants of the case system left than English and we will even devote lesson 15 to those remnants. The personal pronouns actually still show some case differences in both languages. Personal pronouns are quite familiar in English: They are words like I,you,he,she,we,you and they. &lt;br&gt;At least this is the case for the subject (nominative case). As object (accusative) some of them are different: "me",you,"him","us",you,"them". Compare: Notice how I turns into me when used as an object. You remains the same. Much like in English, ik (subject) turns into mij as object in Dutch, whereas je remains the same in both roles: The system in Dutch resembles the English one quite a bit, after all the languages are close relatives: Nevertheless the Dutch system is a little more involved, as we have seen there are: In addition there are Let's leave most of these complications aside for the moment and concentrate on the forms of the pronouns. There are a few more than in English. Exercises 1-1. Quizlet. This is the point where it is your turn to put in some effort yourself, because obviously you have some memorization to do. There is a Quizlet practice set (27 terms) to help you with memorizing the pronouns. But it is recommended to "first" use the above tables. Unfortunately, the pronunciation of some of the clitics with apostrophes is wrong at Quizlet. So click the arrow buttons here to listen to the pronunciation and speak it out loud yourself until you feel confident that you know them, then go to Quizlet. Do make sure you can hear sound at Quizlet.First scan through the cue cards, then do some of the other methods available. It should take you an hour or so and you will know some of the most frequently used words in the language. Woordenlijst 1. You have already encountered quite a few words above. Now make sure you own them! Listen to their pronunciation, sort the table by English and read back to Dutch, check the pronunciation again. Click on the blue link to go to the Dutch wiktionary and try to figure out what you may. If you do not understand, follow the interwiki link to go to the English wiktionary. In short: there are many ways to use this table and you can try one thing one day and come back another to try something different. &lt;br clear="all"&gt; Quizlet. The vocabulary of this lesson can be trained at Quizlet. (33 terms) Your turn! Building vocabulary 1. When learning a language you need to start building up your vocabulary. There are various ways of doing that. One is to study the above conversations well. Often words are easier to remember when put in context. We will add vocabulary building exercises to each lesson to make it easier for you to memorize it all. Progress made. If you have studied the above well, you should Further practice. This lesson is accompanied by two pages that are intended to practice and reinforce what you have learned above. They do that with a bit different approach It is recommended that you first work on the material in these two modules before you move on to lesson two, but of course this depends on your level of understanding and one of the nice things about the wiki-system is that one can use it whichever way you see fit. (Which is what children would do, but they are used to running into new things that they do not fully comprehend.) 

Nursery rhymes. Children learn a lot of language skills by playing, singing, dancing. They know how to make learning fun. This is why children's songs and rhymes are a wonderful way to acquire a foreign language. Here are three examples. Enjoy being a child again! Poesje en Hondje. The following text was taken from a Mother Goose rhyme and translated into Dutch. In order to get a literal translation, the Dutch text was not made to rhyme. Note that in Dutch the word "poesje" does "not" have the same connotation as in English. It merely means "pussycat". Poesje Mauw. The following is a Dutch "volksliedje" (folk song). &lt;br clear="all"&gt; The meter in the Dutch version is nearly perfect and should provide hints for pronouncing the words. Lekkere is pronounced as 'le-kre' in this case, to fit the meter (but this poetic license, not non-standard pronunciation). Now that you understand the poem, go see a video of it, see here (Notice that in some dialects "ij" and "ei" are pronounced more like [ɑɪ̯] than as [ɛɪ̯].) There is a pretty astounding 'performance' of this song by (the late) Corrina Konijnenburg that was recorded in 1967 in a children's show by "Dorus" (real name Tom Manders). Note that the performer took a few liberties. Notice that she pronounces poesje as "poessie" as is usual in Hollandic dialects. Ba, ba, black sheep. This nursery rhyme is well known in English, but here is the Dutch version. &lt;br clear="all"&gt; Quizlet. The vocabulary of this lesson can be practiced at Quizlet (21 terms) Progress made. If you have studied the above well you should have Cumulative count: Les 1: 116 terms, Les 1A: 89 terms. Example 1: 21 terms Total 226 terms. Further learning. Please proceed to Dutch/Lesson 2 

A route is the path that data takes when travelling through a network from one host to another. Routing is the process by which the path, or some subset of it, is determined. One of the characteristic features of the Internet, as compared to other network architectures, is that each node that receives a packet will typically determine for itself what the next step in the path should be. IP routing decisions are generally made based on the destination of network traffic. When an IP packet is sent from a node on the network, it will consult its "routing table" to determine the next hop device that the traffic should be sent to, in order for it to reach its final destination. The routing table on a typical home machine may look something like this (except formatted properly :):  Kernel IP routing table  Destination Gateway Genmask Flags Metric Ref Use Iface  x.y.z * 255.255.255.255 UH 0 0 0 ppp0  192.168.0.0 * 255.255.255.0 U 0 0 0 eth0  127.0.0.0 * 255.0.0.0 U 0 0 0 lo  default x.y.z 0.0.0.0 UG 0 0 0 ppp0 So, for example, when it receives a packet on interface eth0 which has a destination of 216.239.59.104, it will consult the table and see that it should send it through the default interface, the host x.y.z, which is on interface ppp0. The routing table is constructed from a combination of statically defined routes and those learned from dynamic routing protocols. Statically defined routes may be declared at system boot time, or via a command line interface. They will generally include the following parameters: Static routes may also include the following parameters: The "default route" is a special case of a statically defined route. It is the route of last resort. All traffic that does not match another destination in the routing table is forwarded to the default gateway. Dynamic routing protocols allow network attached devices to learn about the structure of the network dynamically from peer devices. This reduces the administrative effort required to implement and change routing throughout a network. Some examples of dynamic routing protocols are: ISIS and OSPF are link-state protocols, meaning each node part of the same zone, will know the state of all the link in the mesh. Due to the exponential number of link in a mesh, thoses protocols are for small mesh such as an ISP national backbone. RIP is usually used to easily announce customer's routes in a backbone. BGP is used as an external routing protocol to exchange routes with other entities. ISP use BGP extensively to "trade" their routes. It can also be used to carry customers routes accross a network, in a MPLS backbone for example. 



Articles Articles can be definite or indefinite. In both cases, as in the other Romance languages of the Iberian Peninsula, articles have two genders: masculine and feminine. Number can be singular or plural. The most general paradigm for the definite article is: Gender and number must agree with that of the noun (and adjective). 

 Lektion 7  Einfache Mathematik ~ Simple Mathematics Lernen 7 ~ Zählen von 13 bis 100. Once you have memorized the numbers from 1 to 12 (see Lernen 3), counting higher in German becomes very much like counting in English. From 13 to 19, add "-zehn" (10; "-teen" in English) after the cardinal number root: 13 – "dreizehn" (irregular in English: 'thirteen')&lt;br&gt; 14 – "vierzehn"&lt;br&gt; 15 – "fünfzehn"&lt;br&gt; 16 – "sechzehn" (note that the 's' in "sechs" is dropped and the 'ch' is pronounced like the 'ch' in "ich")&lt;br&gt; 17 – "siebzehn" (note that the 'en' in "sieben" is dropped)&lt;br&gt; 18 – "achtzehn"&lt;br&gt; 19 – "neunzehn" Above 19 the counting system is constant: add "-zig" ("-ty" in English) to the cardinal root. Thus, we get: 20 – "zwanzig"&lt;br&gt; 21 – "einundzwanzig" (note: 'one-and-twenty')&lt;br&gt; 22 – "zweiundzwanzig" (note: 'two-and-twenty')&lt;br&gt; And the same for 30, 40, 50...etc. 30 – "dreißig" (this is an exception to the -zig Rule)&lt;br&gt; 40 – "vierzig"&lt;br&gt; 50 – "fünfzig"&lt;br&gt; 60 – "sechzig"&lt;br&gt; 70 – "siebzig"&lt;br&gt; 80 – "achtzig"&lt;br&gt; 90 – "neunzig"&lt;br&gt; 100 – "hundert" So, combining these, we get: 34  – "vierunddreißig" (note: 'four-and-thirty')&lt;br&gt; 143 – "hundertdreiundvierzig" (note: 'hundred-three-and-forty')&lt;br&gt; 170 – "hundertsiebzig"&lt;br&gt; 199 – "hundertneunundneunzig" It would be excellent practice towards learning these numbers by counting (in German, of course) from 1 to 199—or counting along any continuous sequence that comes to mind. For example, start with your age and count to 50 (count down if appropriate). Grammatik 7-1 ~ Math Calculations. The following table presents the symbols used for basic mathematics. We can use these symbols to ask and answer simple problems in mathematics. Some of the examples that follow include first a question ("Frage") and then the answer ("Antwort"):  "Wieviel ist sechs und sieben?" How much is 6 and 7?  "Sechs und sieben ist dreizehn" 6 and 7 is 13  "Wieviel ist fünfzig plus achtzehn?" How much is 50 + 18?  "Fünfzig plus achtzehn ist gleich achtundsechzig" 50 + 18 = 68  "Wieviel ist siebzig minus zehn?" How much is 70 - 10?  "Siebzig minus zehn ist gleich sechzig" 70 - 10 = 60  "Wieviel ist neun durch drei?" How much is 9 divided by 3?  "Neun durch drei ist gleich drei" 9 ÷ 3 = 3  "Funf ist größer als zwei" 5 &gt; 2  "Acht ist kleiner als siebzehn" 8 &lt; 17 Vokabeln 7-1. Counting to 199 is also included in the vocabulary for "Lektion 7".  die Antwort answer  die Frage question  geteilt/dividiert durch over [math]  größer als greater than  kleiner als smaller than  geteilt/dividiert divided, forked, split  gleich equal, same, even  hoch tall, to the power of [math]  mal times [math]  minus minus  plus plus  wieviel? how much? 

Hub Bearings. Over time the bearings in the hubs of a bicycle wheel may come out of adjustment causing the wheel to wobble from side to side (this can also be the result of the wheel being out of true; a sure sign of loose bearings is that the rim can be moved laterally in the fork by hand). Additionally, road dirt and moisture infiltrate the bearings, causing rough operation and premature wear. Even if these issues do not arise, the bearings' lubrication will eventually need to be replaced in order to maintain the life and health of the hubs. These problems can be addressed by overhauling the hubs. The basic techniques are similar to maintaining any other ball bearing assemblies, whether the headset or bottom bracket on the bike, or on complete different applications. On some recent model mountain bikes, there has been a rash of rear hubs getting significant play after only a few rides; this may be due to bad cones or lock nuts. The drive side lock nut should be your number one suspect in this case if the bike is new. Parts of the hub. From the inside out: Shell. Main body of hub, holds the axle assembly and is the connection point for the spokes. Bearing Cup. Is pressed into the shell. Bearing Cone. Forms the outside bearing races, fits onto axle, usually adjustable. Bearings. There are three major types of bearings in use on bicycle hubs: Cleaning and repacking a traditional cone and cup front hub. At this point, depending on how the grease is distributed on the parts, the ball bearings may well stick to the axle, or drop out of the hub or into the axle hole through the hub, or remain in the cup. If the bearings are in a cage, they will generally remain in place unless the cage has been excessively worn. It is advisable to leave the other cone and locknut assembly undisturbed in its position on the axle, as this will allow you to maintain the original side to side alignment of the axle assembly in the hub as much as possible. (Of course, this is mandatory where this assembly is not removable). It can be cleaned and greased in this condition. (If the cone is damaged, then it will need to be disassembled and replaced). Maintaining a cartridge bearing hub. There are many different designs so this is just an overview of the basic principles that can be used to service most cartridge bearing hubs. Clean and lubricate bearings. If you ride in wet conditions you may be able to extend the life of the bearings by periodically cleaning and lubricating them. In most circumstances this is not required and the bearing will have an acceptable life with no maintenance. The idea is to remove one of the seals from the bearing so it can be cleaned and fresh grease applied. Replace bearings. Replacement bearings can be sourced from good bike shop or from any bearing supplier. Typically bearings are identified from the numbers on the side of the bearing, take the old bearings with you to make sure you get the correct replacements. Removing old bearings. Installing new bearings. Cartridge bearings are not good at absorbing side loads. When being pressed into the hub the force needs to be pushing on the outer metal ring to avoid damaging the bearing. A suitable sized socket can be used to do this. 

Cleaning parts can be done with a cleaning solution of some kind, and/or with a rag or stiff brush. A couple of notes about cleaning products: Power Washer. Using a power washer is an option in some situations, but care must be taken not to displace lubricant from parts. A high pressure stream of water is capable of displacing the grease from a bearing surface, including from sealed bearings, which are not designed to be lubricated by the user. Some very powerful power washers can also take the paint off of a bike. As such it is not recommended to clean a bike with a power washer because it is very hard to dry the water and replace the grease. Cleaning and lubricating a bicycle chain. Cleaning and lubricating your chain is one of the simplest and most important things to do to keep your bike in good working order. It is also an activity that draws intense debate with many people insisting that their way is the "Right" way to do it. The short story is, the cleaner and better lubricated your chain is the longer the entire drivetrain (the chain and gears) will last, but no drive-train lasts forever. It is worth spending time taking care of your drivetrain, it's not worth obsessing about. You need: The choice of lubricant is important. DO NOT use WD40* as it is too thin and will evaporate leaving nothing to lubricate the chain. Use special bike chain lube found at any bike store. Some lubricants are formulated for wet conditions, others for dry. The difference is the viscosity. A dry lube used in wet conditions will require more frequent applications, a wet lube used in dry conditions will gather dirt, and require more frequent cleaning. Prior to lubricating, clean the chain as well as possible. Usually a chain cleaning tool, or a toothbrush is sufficient. Remove the chain from the bicycle and soak it in a solvent such as a cycle degreaser. Hang up chain to dry off. Paraffin will also work, but is unpleasant and unhealthy to work with environmentally hazardous, and flammable; as such it is not recommended. Dirt is what causes a chain to wear out, so keeping your chain as clean as possible extends its life, and the life of the entire drivetrain. Apply the lubricant on the chain non-pressurized drip type lubes are recommended as are more accurate, and waste less lube than spray lubricants. Put a rag behind the chain so that the lube does not contaminate the tires or other parts of the bike. Rotate the pedals to move the chain and lubricate the whole length. Wipe the chain with the rag to remove the excessive lubricant, as this just gathers dirt and does not help lubricate the chain. 

 Lektion N Richtig oder falsch? 1 a Wir müssen gehen. So wir haben keine Zeit. 1 b Wir müssen gehen. Deshalb haben wir keine Zeit. 2 a Wenn wir haben Stress, machen wir oft Fehler. 2 b Wenn wir Stress haben, machen wir oft Fehler. 3 a Ich mag Leute nicht, die zu spät kommen. 3 b Ich mag Leute nicht, die kommen zu spät. Richtig sind: 1 b, 2 b, 3 a, 

 Lektion NA &lt;/small&gt; "Contributed to Lesson 1A, but concept not really introduced until Lesson 5":  Beispiel an "können" (can) "gehen" (go) "sein" (to be)  pronoun verb I (irreg?) verb II verb III (irregular)  Basicform können gehen sein  ich kann gehe bin  du kannst gehst bist  er/sie/es kann geht ist  wir können gehen sind  ihr könnt geht seid  sie können gehen sind  Sie (formal) können gehen sind As you can see, any verb uses the same declination for wir, sie and Sie. Also, er/sie/es uses the same declination for all three genders. NOTE: Moved to Lesson 6 "Used this story for German/Level III/Markus Studiert since it nicely followed your other Markus conversation. But I could not match "holt" or "holt aus" to any verb. Can you help? " -- holt is from holen, holt aus is from ausholen, which has another meaning, but "holen" is meant here in the sense of "fetching a book out of the shelf" (trying to match the German sentence literally) Markus ist in der Universität. Er trinkt dort einen Kaffee und ißt ein Brötchen. Danach geht er in die Bibliothek. Er sucht ein Buch über Biochemie. Er holt das Buch aus dem Regal und setzt sich an einen Tisch. Nach einer Stunde geht er in den Hof und raucht eine Zigarette. Danach geht er an den Tisch zurück. Er denkt: "Wenigstens eine Stunde..." und stellt das Buch wieder in das Regal. Heute ist nicht sein Tag. "There are already lessons completed for some of these points. However, it might be that the stories or conversations need to made easier to read in those lessons." What should be explained here? "Note that the book layout is such that the Advanced lessons add to the basic lesson (1A furthers something taught in 1), but 1A cannot be the place where anything basic is introduced, as the berginning student should not even visit Lesson 1A until later". Just wondering, but isn't this very limited for Advanced Level? I would like like some help for the Wirtschaftsdeutsch Prüfung but this isn't going to help me. 

Hi there! Welcome to this Wikibook on the wonderful world of Bourne Shell Scripting! This book will cover the practical aspects of using and interacting with the Bourne Shell, the root of all shells in use in the world. That includes interacting with the shell on a day-to-day basis for the purposes of operating the computer in normal tasks, as well as grouping together commands in files (scripts) which can be run over and over again. Since it's not practical to talk about the Bourne Shell in complete isolation, this will also mean some short jaunts into the wondrous world of Unix; not far, just enough to understand what is going on and be able to make full use of the shell's very wide capabilities. There are also some things this book won't do for you. This book is not an in-depth tutorial on any kind of programming theory -- you won't learn the finer points of program construction and derivation or the mathematical backings of program development here. This book also won't teach you about or any other type of Unix or Unix itself or any other operating system any more than is necessary to teach you how to use the shell. Nothing to be found here about , joe, vi, or any other specific program. Nor will we cover firewalls and networking. We "will" cover the Bourne Shell, beginning with the basic functionality and capabilities as they existed in the initial release, through to the added functionality specified by the international POSIX standard POSIX 1003.1 for this shell. We will have to give you "some" programming knowledge, but we hope that everyone will readily understand the few simple concepts we explain. Having said that, the authors hope you will find this book a valuable resource for learning to use the shell and for using the shell on a regular basis. And that you might even have some fun along the way. 

This is a summary of Scott McCloud's "Understanding Comics". Chapter 1: Setting the Record Straight. Comics is a medium, not a genre. In one of the few previous books discussing comics as a medium, Will Eisner's "Comics and Sequential Art", comics is defined, unexpectedly, as Sequential Art. Here McCloud expands and formalizes that definition (in a rare panel that reduces neatly to pure text): "com.ics (kom'iks) n. plural in form, used with a singular verb. 1. Juxtaposed pictoral and other images in deliberate sequence, intended to convey information and/or to produce an aesthetic response in the viewer." The definition is dense (in the book it is developed over several pages, with McCloud's cartoon avatar taking questions from an audience), but to quote another part of the chapter1 The secret is not in what the definition says but what it "doesn't" say! || For example, our definition says nothing about superheroes or funny animals. Nothing about fantasy/science fiction or reader age. No genres are listed in our definition, no types of subject matter, no styles of prose or poetry. || Nothing is said about paper and ink. No printing process is mentioned. Printing "itself" isn't even specified!2 Nothing is said about technical pens or bristol board or Windsor &amp; Newton Finest Sable Series 7 Number 2 Brushes!3 No materials are ruled out by our definition. No tools are prohibited. || There is no mention of black lines and flat colored ink. No calls for exaggerated anatomy or for representational art of any kind. No schools of art are banished by our definition, no philosophies, no movements, no ways of seeing are out of bounds!4 Like text, comics can be used in uncountable ways; unlike text, its potential has until recently been tragically squandered. Chapter 2: The Vocabulary of Comics. Comics are built out of icons -- "For the purposes of this chapter, I'm using the word icon to mean any image used to represent a person, place, thing or idea." Letters and words are completely abstract icons, bearing no physical resemblance to their ideas; pictures have varying levels of abstraction, from the photograph of a face, full of color and shading, an ink copy of photons taken in through a shutter, to =D That smiley, that :-) , that set of dots and lines that looks like a face to us only because our brains contain a lot of face-seeing hardware (after all, no one sees eyes in solitary colons, noses in solitary dashes), is a cartoon. Cartoons focus our attention -- through simplification, by eliminating superfluous features, they amplify the features that remain -- but, McCloud explains, that is not the entirety of their drawing power: When two people interact, they usually look directly at one another, seeing their partner's features in vivid detail. || Each one also sustains a constant awareness of his or her own face, but "this" mind-picture is not nearly so vivid; just a sketchy arrangement...a sense of shape...a sense of "general placement". || Something as simple and basic as a cartoon. || Thus, when you look at a photo or a realistic drawing of a face, you see it as the face of another. || But when you enter the world of the "cartoon", you see "yourself". ... The cartoon is a vacuum into which our identity and awareness are pulled, || an empty shell that we inhabit which enables us to travel into another realm. We don't just observe the cartoon, we "become" it! The spectrum as described thus far: real face --&gt; photograph of a face --&gt; realistic drawing --&gt; cartoony drawing --&gt; smiley The photograph requires little cognition; it is received as pure visual input. The smiley requires an effort, if an unconscious one, a piecing together of the abstract lines. Continuing the spectrum to the right, McCloud argues, we arrive at the word FACE -- bold and distinct, with relatively few letters, reminiscent of prehistoric times, when words and pictures were one -- then at a more detailed description ("Two eyes, one nose, one mouth" is the example given), then at still higher levels of abstraction ("Thy youth's proud livery, as gazed on now..."). reality --&gt; photo --&gt; Batman --&gt; Charlie Brown --&gt; =D --&gt; FACE --&gt; Two eyes, one nose, one mouth --&gt; "Thy youth's proud livery, as gazed on now..." Our need for a unified "language" of comics sends us toward the center where words and pictures are like one side of the same coin! || But our need for "sophistication" in comics seems to lead us outwards, where words and pictures are most separate. In addition to iconic abstraction, there exists abstraction of the more traditional sense, the abstraction of abstract art. Thus, the spectrum becomes a pyramid:  The Picture Plane/Art Object  / .\ \  / \ \  / . \ .|  / . \ / language  reality The pyramid has limitations, of course (Where do The Treachery Of Images and Fountain fall? Does the horizontal axis along the verbal edge dictate how the words are displayed on the page or how they describe things?) but it is a cool little invention. Chapter 3: Blood in the Gutter. As part of normal life, everyone learns to assume certain things (or so one assumes) -- that the world doesn't disappear when you're not looking, for example, that the house across the street has furniture and interior walls. "As infants, we're unable to commit that act of faith. If we can't see it, hear it, smell it, taste it or touch it, it isn't there! The game "Peek-A-Boo" plays on this idea. Gradually, we learn that even though the sight of mommy comes and goes, mommy remains. || This phenomenon of observing the parts and perceiving the whole has a name. It's called closure." You performed closure when you saw the lines at the top of this section as two anime smilies; more to the point of the chapter, you performed closure when you saw the two smileys as a single "winking" smiley. "See that space between the panels? That's what comics aficionados have named "the gutter!" And despite its unceremonious title, the gutter plays host to much of the magic and mystery that are at the very heart of comics...If visual iconography is the vocabulary of comics, closure is its grammar." Of course, other media make use of closure as well -- in movies, our minds effortlessly connect each frame to those preceding and following it -- but comics requires conscious (or semiconscious), high-level closure between every frame. McCloud has categorized panel-to-panel transitions into six classes: Looking at how often each panel transition is used in a particular comic can reveal some interesting things. Jack Kirby's pioneering style, as invoked in a Fantastic Four comic from 1966, breaks down as follows: 65% action-to-action (type 2), 20% subject-to-subject (type 3), 15% scene-to-scene (type 4); the remaining transitions are unused. Here's the bar graph McCloud makes of the data:  |1|2|3|4|5|6|  90%| | | | | | |  | |.| | | | |  | |M| | | | |  50%| |M| | | | |  | |M| | | | |  | |M| | | | |  | |M|M|.| | |  10%| |M|M|M| | | As it turns out, almost every American comic -- regardless of storytelling style, regardless of genre (with a few experimental exceptions like Art Spiegelman's early work) -- charts similarly, from issue 1 of X-Men to Heartbreak Soup to Betty and Veronica to Naughty Bits to Frank in the River to A Contract With God to Maus to Donald Duck. A similar survey of European comics (Squeak the Mouse, Asterix, Welcome to Alflolol, The Long Tomorrow, Manhattan, Clik!, The Black Island, The Clock Strikes) yields similar results. But here's a popular mainstream Japanese comic from Osamu Tezuka:  |1|2|3|4|5|6|  90%| | | | | | |  | | | | | | |  50%| | | | | | |  | |M| | | | |  | |M|M| | | |  | |M|M| |L| |  10%|L|M|M|M|M| | In Japan, where comics developed mostly in isolation following World War II, where they are often published in gigantic anthologies rather than tiny magazines (lessening the premium on space and thus the emphasis on concise, action-oriented transitions 2-4) the charts look quite different. More important still, eastern culture has bequeathed an emphasis on holism and contemplation (aspect-to-aspect works well for setting a mood), an emphasis on the power of intervals -- of silence in a song, of negative space in a painting. In comics this means a renewed emphasis on the power of closure, on the strange alchemy that occurs in the gutter. The effect is to spark the imagination, to engage every one of the five senses, rather than simply sight. (Sadly, replicating McCloud's demonstration of this in ASCII art is prohibitively difficult.) The comics creator asks us to join in a silent dance of the seen and unseen. The visible and invisible. || This dance is unique to comics. No other artform gives so much to its audience while asking so much from them as well. || This is why I think it's a mistake to see comics as a mere hybrid of graphic arts and prose fiction. What happens between these panels is a kind of magic only comics can create. Chapter 4: Time Frames. We've been trained to see a picture as a snapshot, as a single moment in time, but any student of visual physiology can tell you that the vast majority of information is taken in by a small area at the center of the field of vision; our eyes compensate for this porthole effect by darting continually around and our brains compensate by maintaining the illusion of detailed peripheral vision (for a demonstration, try reading this sentence while keeping your eyes fixed on one of its component words). Much modern art acknowledges this (, for example, incorporates many perspectives into one image), and so do many comics. McCloud's example is a long panel containing several chatting family members; reading from left to right, the conversation takes maybe 15 seconds to unfold, and the expression and pose of each person matches the moment his or her particular statement is made -- "one panel, containing several panels". To some extent, then, space in comics translates into time; as your eyes cross the page they also pass through the seconds (or the hours, or the years) and while a frame generally denotes a particular moment (in addition to serving other purposes beyond the scope of this entry), "moment" is a slippery thing that in practice can be any length at all; in determining a more specific time period, the reader relies heavily on context, on a page's particular content, and, especially, on the sounds portrayed in the text, which we have not been conditioned to think of as having a duration of zero. And then there's motion. It can, of course, be portrayed through the frame transitions described in chapter 3, but shortly after the era of futurism had ended, shortly after the invention of the motion picture, comics invented the motion line, which lies "somewhere between the futurists' "dynamic" movement and Duchamp's diagrammatic "concept" of movement." Over a period of decades, motion lines evolved from "wild, messy, almost desperate attempts to chart the paths of moving objects through space" to something "more refined and stylized, even "diagrammatic"", then, eventually, "became "so" stylized as to almost have a life and physical presence all their own!" I've been trying to figure out what makes comics tick for years and I'm still amazed at the strangeness of it all. || But no matter how bizarre the workings of time in comics is -- || -- the face it presents to the reader -- || -- is one of simple normalcy. || Or the illusion of it, anyway. || It all depends on your frame of mind. Chapter 5: Living in Line. This chapter relies very heavily on visual examples, so is difficult to summarize here. Chapter 6: Show and Tell. In the beginning, pictures and words were two sides of the same coin. In school, at show-and-tell, you show, and you tell, interchangeably -- "This is my robot...it's got one of "these" things"; in picture books, simple words combine similarly with images. As we grow up, we learn to separate "show" and "tell" from each other -- we paint pictures without words, and read books without pictures. A similar progression can be traced through history. In cave paintings, the people shown were iconic, symbol-esqe, almost like letters, as were the flat, bright images of the early Greeks and Creteans and the line-drawings of ancient Egypt; early written language, likewise, was full of the descendants of those cave paintings, letters that were pictures. Words and images were side by side, at the lower-left vertex of McCloud's great pyramid. Over the course of the next few thousand years, they diverged. Letters sacrificed visual representation for writing ease (and, later, printing ease) and pictures grew richer and more complex until looking at them was more like looking at reality than at thoughts. 

Applications of Special Relativity. In this chapter we continue the study of special relativity by applying the ideas developed in the previous chapter to the study of waves. First, we shall show how to describe waves in the context of spacetime. We then see how waves which have no preferred reference frame (such as that of a medium supporting them) are constrained by special relativity to have a dispersion relation of a particular form. This dispersion relation turns out to be that of the relativistic matter waves of quantum mechanics. Second, we shall investigate the Doppler shift phenomenon, in which the frequency of a wave takes on different values in different coordinate systems. Third, we shall show how to add velocities in a relativistically consistent manner. This will also prove useful when we come to discuss particle behaviour in special relativity. A new mathematical idea will be presented in the context of relativistic waves, namely the spacetime vector or four-vector. Writing the laws of physics totally in terms of relativistic scalars and four-vectors ensures that they will be valid in all inertial reference frames. Waves in Spacetime. Waves in Spacetime We now look at the characteristics of waves in spacetime. Recall that a wave in one space dimension can be represented by formula_1 where formula_2 is the (constant) amplitude of the wave, formula_3 is the wavenumber, and formula_4 is the angular frequency, and that the quantity formula_5 is called the phase of the wave. For a wave in three space dimensions, the wave is represented in a similar way, formula_6 where formula_7 is now the position vector and formula_8 is the wave vector. The magnitude of the wave vector, formula_9 is just the wavenumber of the wave and the direction of this vector indicates the direction the wave is moving. The phase of the wave in this case is formula_10. &lt;br&gt;Figure 5.1: Sketch of wave fronts for a wave in spacetime. The large arrow is the associated wave four-vector, which has slope formula_11. The slope of the wave fronts is the inverse, formula_12. In the one-dimensional case formula_5. A wave front has constant phase formula_14, so solving this equation for formula_15 and multiplying by formula_16, the speed of light in a vacuum, gives us an equation for the world line of a wave front: formula_17 The slope of the world line in a spacetime diagram is the coefficient of formula_18, or formula_19, where formula_20 is the phase speed. 

Grammatica 2-1 ~ Introduction to Verbs. A (in Dutch: "werkwoord") is that part of speech that describes an action. Verbs come in an almost bewildering array of tenses, moods, voices and aspects, and there are several major types: intransitive, transitive, ditransitive, and ergative verbs. Fortunately, the Dutch verb is not too different from the English one, although it does have a few more forms.  I am called Standish "Ik heet Standish"  What are you called (named)? "Hoe heet u?"  ...that she is named (called) 'Alice' "...dat ze 'Alice' heet"  We are both called Robert "Wij heten allebei Robert" The Dutch verb "heten" can best be translated as "to be named" or "to be called" and we see two forms of it here Actually there are usually three forms. This can be seen from:  In the case of "heten" the extra -t does not get added because the stem already ends in a "-t". In a later lesson we will revisit the verb forms associated with each person. The irregular verb "to be"-"zijn" has a few more forms in both languages. Gesprek 2-2 ~ De Engelsman. First push the arrow button to listen to the following conversation. Then inspect the translation and hover over each word you do not know to find out what it means. Once you understand the narrative run the audio again, following along, making sure you know what is being said. Use the pronunciation box on the right to further strengthen your comprehension both in listening and in reading. Fill in the blank 2-2-F. Say the word you think that belongs in the blank and use the hover method to check your choice Vocabulary drill 2-1. Of course memorizing words and expressions is an important part of learning any language and there are various ways of doing that. Have a look at the vocabulary pages. They are designed to help you acquire more words in a playful manner. Grammatica 2-2 ~ Inversion in questions and negations. You may have wondered about the order of the words in Even though Dutch verbs are not so much more complicated than English ones, word order is. In fact it is quite a bit more complicated than in English. For the moment let's just leave the above sentence for what it is and start with questions. Questions. A question sentence in Dutch simply reverses the order of subject and verb. Recall: "U heet meneer Standish" ('You are named Mr. Standish). It became: "Hoe heet u?" as a question The normal word order of subject ("u" or "you") then verb ("heten") is "reversed" and, in this case, an interrogative ("hoe" or "how") added. Additional examples: English does the same thing when using the verb "to be": Dutch does not use the auxiliary "to do" as English requires in most other cases: Negations. The negative is formed by simply adding "niet" at the end: Again, Dutch does "not" use the auxiliary "to do". (In fact using it sounds very foreign.) Even a negative question does not use "to do": Gesprek 2-3 ~ Het nieuwe meisje. In this conversation, the parties are close friends. Fill in the blank 2-3-F. Use the hover method to check your answer. Grammatica 2-2 Adjectives, demonstratives and articles. Gender. Where English uses the demonstrative pronoun "that", Dutch uses either "dat" or "die", recall: Similarly, where English uses the article "the", Dutch has two possibilities: "de" or "het", recall: We will revisit this phenomenon (gender) in the next lesson more extensively. There is a bit of a problem with it in Dutch. For the moment it is enough to realize that there are two kinds of words, For this reason it is advisable to always memorize a word "together" with its definite article, e.g. as "de boekhouding", not simply as "boekhouding". Both articles and demonstrative pronouns are a special kind of adjectives: words that are added to make the meaning of another word more precise, like "new", "small" or "exciting" Inflection. Recall that some adjectives in the dialogue ended in -e (mooi"e" meid), sometimes they did not (is erg mooi). Adjectives can be used in two ways: in front of a noun and after a verb like "is" (a copula). In English the adjective remains the same regardless: Behind a copula (as "predicate") this is true in Dutch as well: But in Dutch they are "inflected" if they occur in front of a noun (as "attribute"). Compare: Neuter words are the ones that carry the definite article "het" and the demonstrative "dat". They are a bit different (Again: we will revisit them in the next lesson.) As you see the adjective is not inflected after the indefinite article "een". This also holds if there is no article: But: Thus, apart from the indefinite neuter an attributive adjective is usually inflected with -e. There are a few exceptions, compare e.g.: Making nouns out of adjectives. Adjectives can be turned into nouns, by assuming their inflected form: Notice that Dutch does not use 'one' in such cases. There are a number of adjectives that can be turned into nouns by adding -te. They all carry de. In English the corresponding suffix is -th: More about nouns in the next lesson. Woordenlijst 2. &lt;br clear="all"&gt; Quizlet. The vocabulary can be practiced as Quizlet (30 terms) Progress made. If you have studied the above lesson well you should have Cumulative vocabulary count: 

Almost all books like this one have a section on (or very similar to) "why you should use the shell/program flavor/language/etc. discussed in this book and not any of the others that perform the same tasks in a slightly different way". It seems to be pretty well mandatory. However, this book will not do that. We'll talk a bit about "why Bourne Shell" of course. But you'll soon see that doesn't preclude other shells at all. And there's no good reason not to use another shell either, as we will explain a little further down. Bourne shell and other Unix command shells. There are many Unix command shells available today. Bourne Shell is just one drop in a very large ocean. How do all these shells relate? Do they do the same things? Is one better than the other? Let's take a look at what makes a shell and what that means for all the different shells out there. How it all got started.... The Unix operating system has had a unique outlook on the world ever since it was created back in the 1970s. It stands apart from most other operating systems in that its focus has always been towards "power users": people who want to squeeze every drop of performance out of their system and have the technical knowledge to do so. Unix was designed to be programmed and modified to the desires of the user. At its core, Unix does not have a user interface; instead it is comprised of a stable OS kernel and a versatile C library. If you're not trying to do actual hard-core programming but rather are trying to do day-to-day tasks (or even just want to put a little program together quickly), pure Unix is a tremendous pain in the backside. In other words, it was clear from the start that a tool would be needed that would allow a user to make use of the functions offered him by the coding library and kernel without actually doing serious programming work. A tool in other words that could pass small commands on to the lower-level system quickly and easily. Without the need for a compiler or other fancy editor, but with the ability to tap into the enormous power of the underlying system. Stephen Bourne set himself to the task and came up with what he called a "shell": a small, on-the-fly compiler that could take one command at a time, translate it into the sequence of bits understood by the machine and have that command carried out. We now call this type of program an interpreter, but at the time, the term "shell" was much more common (since it was a shell over the underlying system for the user). Stephen's shell was slim, fast, and though a bit unwieldy at times, its power is still the envy of many current operating system command-line interfaces today. Since it was designed by Stephen Bourne, this shell is called the Bourne Shell. The executable is simply called "sh" and use of this shell in scripting is still so ubiquitous, there isn't a Unix-based system on this earth that doesn't offer a shell whose executable can be reached under the name sh. ...And how it ended up. Of course, everyone's a critic. The Bourne Shell saw tremendous use (indeed, it still does) and as a result, it became the de facto standard among Unix shells. But all sorts of people almost immediately (as well as with use) wanted new features in the shell, or a more familiar way of expressing commands, or something else. Many people built new shells that they felt continued where Bourne Shell ended. Some were completely compatible with Bourne Shell, others were less so. Some became famous, others flopped. But pretty much all of them look fondly upon Bourne Shell, the shell they call "Dad..." A number of these shells can be run in sh-like mode, to more closely emulate that very first sh, though most people tend just to run their shells in the default mode, which provides more power than the minimum sh. It's Bourne Shell, but not as we know it.... So there are a lot of shells around but you can find Bourne Shell everywhere, right? Good old "sh", just sitting there faithfully until the end of time... Well, no, not really. Most of the sh executables out there nowadays aren't really the Bourne Shell anymore. Through a bit of Unix magic called a link (which allows one file to masquerade as another) the sh executable you find on any Unix system is likely actually to be one of the shells that is based on the Bourne shell. One of the most frequently used shells nowadays (with the ascent of free and open-source operating systems like GNU and Linux) is a heavily extended form of the Bourne Shell produced by the Free Software Foundation, called Bash. Bash hasn't forgotten its roots, though: it stands for the Bourne Again SHell. Another example of a descendant shell standing in for its ancestor is the Korn Shell (ksh). Also an extension shell, it is completely compatible with sh -- it simply adds some features. Much the same is true for zsh. Finally, a slightly different category is formed by the C Shell (csh) and its descendant tcsh, native on BSD systems. These shells do break compatibility to some extent, using different syntax for many commands. Systems that use these shells as standard shells often provide a real Bourne Shell executable to run generic Bourne Shell scripts. Having read the above, you will understand why this book doesn't have to convince you to use Bourne Shell instead of any other shell: in most cases, there's no noticeable difference. Bourne Shell and its legacy have become so ingrained in the heart and soul of the Unix environment that you are using Bourne Shell when you are using practically any shell available to you. Why Bourne Shell. So only one real question remains: now that you find yourself on your own, cozy slice of a Unix system, with your own shell and all its capabilities, is there any real reason to use Bourne Shell rather than using the whole range of your shell's capabilities? Well, it depends. Probably, there isn't. For the most part of course, you "are" using Bourne Shell by using the whole potential of your shell -- your shell is probably "that" similar to the Bourne Shell. But there is one thing you might want to keep in mind: someday, you might want to write a script that you might want to pass around to other people. Of course you can write your script using the full range of options that your shell offers you; but then it might not work on another machine with another shell. This is where the role of Bourne Shell as the lingua franca of Unix command shells comes in -- and also where it is useful to know how to write scripts targeted specifically at the Bourne Shell. If you write your scripts for the Bourne Shell and nothing but the Bourne Shell, chances are far better than equal that your script will run straight out of the mail attachment (don't tell me you're still using boxes to ship things -- come on, get with the program) on any command shell out there. 

Lektion Eins für Fortgeschrittene Geschichte 1-3 ~ "Markus studiert". This short story ("Geschichte") is told in the 3rd person (see Grammatik 1-3). Note how this is apparent from both the pronoun ("Er" or "he") and verb forms. Vokabeln 1-3.  die Bibliothek library  die Biochemie biochemistry  das Brötchen roll, biscuit  das Buch book  der Fortgeschrittene advancer  die Fortgeschrittenen advancers (pl.)  die Geschichte story  der Hof courtyard; also court  der Kaffee coffee  die Stunde hour  der Tisch table  das Regal shelf  die Zigarette cigarette  denken think (Er denkt = He thinks)  essen eat (Er isst = He eats)  holen fetch, get (Er holt = He gets/fetches)  rauchen smoke (a cigarette) (Er raucht = He smokes)  sich setzen sit (oneself) down (Er setzt sich = He sits)  stellen place (Er stellt = He places)  suchen seek, search for (Er sucht = He looks for)  trinken drink (Er trinkt = He drinks)  aus out  danach afterwards  dort there  in in  nach after  über about  wenigstens at least, at any rate  wieder again Grammatik 1-3 ~ Personal Pronouns. As in English, personal pronouns exist in three grammatical persons, each with singular and plural number. In Gespräch 1-1 and 1-2, you see only the singular versions. The table here gives also the plural (nominative case only): Grammatik 1-3 ~ Incomplete Sentences. What are we to make of short, incomplete sentences such as that in Gespräch 1-1: 'Und dir?'? This translates as: 'And for you?'. In English and German it is not always necessary to express every part of a sentence, especially in conversation where the words left out are easily understood by both or all parties. Walk up to a stranger and say 'And you?' and a possible response is a hostile 'Out of my face, fool'. But in the conversation between Heinrich and Karl, Heinrich knows that Karl is really meaning: "Und wie geht es dir?", with that part underlined left out of the conversational statement. Note especially that the pronoun 'you' retains its case—its relation to the missing verb from the implied sentence—distinctive in German (that is, "dir" instead of "du") but not so in English (the form "you" covers both cases). Übersetzung 1-2. Although these sentences involve many grammatical concepts that have not been covered, each can be written in German by referring to the example sentences and vocabularies in Lessons 1 and 1A. Using a piece of paper and pencil, translate each of these sentences into German: 

Het alfabet ~ The alphabet. The Dutch alphabet, like English, consists of 26 basic letters. However, there are also a number of letter combinations. The following table includes a listing of all these letters and a guide to their pronunciation. As in English, letter sounds can differ depending upon where within a word the letter occurs. The first pronunciation given below (second column) is that in English of the letter (or combination) itself. Reading down this column and pronouncing the "English" words will recite the alphabet "in het Nederlands" (in Dutch). Note that letter order is exactly the same as in English, but pronunciation is not the same for many of the letters. Trouble areas for Anglophones are marked in red Diacritics. Diacritics are not very numerous in Dutch and they are mostly limited to loans, but the orthography does for example demand a diaeresis, mostly on ë and ï in words like "geërgerd, geïnteresseerd". It marks the boundary of two syllables and is fairly common. It is not to be confused with the German umlaut that only occurs on a few German loans like "überhaupt". Less commonly, the orthography also allows stress marks to be added. Acute accents can be used for that purpose, but only then if: There are some 200 words and word forms that occur in two forms that only differ in stress pattern, where that can be an issue. For diphthongs and doubled vowels both vowels are marked as in "áúto". In principle this holds for the digraph "ij" as well. The "j" should also get an acute, but computers usually do not facilitate that. Here we will write "íj" if the need arises. A stress mark is not be confused with an accent grave, aigu or circonflex as it can occur on the more numerous French loans like "café", "à propos" or "maîtresse". Interestingly Dutch orthography maintains the "î" in the latter, although French dropped it in 1990. The grave is also used on the verb "blèren". It is an indigenous onomatopoeia that mimics the bleating of sheep with an unusual long [ɛ:] sound. The cedille is rare but occurs in a few French and Portuguese loans like "façade, Curaçao" A useful tip is to install a Dutch international keyboard. For Microsoft operating systems this will tell you how to proceed. Nederlandse uitspraak ~ Dutch Pronunciation Guide. Klinkers. Dutch has quite a few vowels (13). To be well understood by a native speaker it is imperative to master them, which can be quite challenge for native speakers of languages that rely more on their many consonants such as Russian. One general observation is that they are always pronounced as more or less "pure" (or only slightly diphthongized) vowels as in French, never quite as drawled or 'chewed upon' as in many varieties of English. Admittedly this does vary from speaker to speaker and region to region, but for Anglophones it is certainly advisable to limit the 'drawling'. Most vowels occur in "pairs" that are traditionally indicated by the terms short and long. Unfortunately, this nomenclature is rather misleading because the difference is not so much a matter of length, but rather a difference in the position of the tongue root (lax vs. tense, or in Dutch: gedekt and open, covered and open). In addition there is a neutral vowel that occurs in almost all unstressed syllables, the schwa /ə/ as is does more or less in English as well. It is spelled with an 'e', so that this letter has three meanings: the above two in stressed syllables, the schwa in unstressed ones. The u. Notice the value of the open letter 'u' in Dutch. As in French it denotes the /y/ sound. Thus, in German it corresponds to ü. It is relatively rare in Dutch because most words that used to have it have shifted it to the diphthong "ui". It occurs mostly at the end of words like "u", "nu" or in front of a w or r: "ruw", "stuur". The oe. This combination represents the [u] sound in German Mut or Spanish tu. Most languages use 'u' as German and Latin. (French uses 'ou', English often used 'oo', although it does have words like "shoe"). The Dutch /u/ sound is strongly rounded and dark with the tongue pretty far retracted back in the mouth. American 'oo' sounds tend to be intermediary between /u/ and /y/ and that is a problem. Compare: The eu. This combination represents another difficult vowel for Anglophones. In German it is written as ö as in Möwe. In IPA as [ø]. Before "r" vowels have a tendency to be modified: The spelling rule. The vowels oe and eu are single, but as we saw above five vowels occur in open/covered pairs: a, e, i, o, u. There is a systematic way in the spelling to indicate which of the two varieties is intended. If a conflict arises, either the vowel or the consonant is doubled: Notice how the formation of the plural necessitates a good mastery of this principle. The vowels oe and eu do not exhibit the dual quality of the other vowels. The case of the letter i is a bit special. There has been a double "ii" in the past but to avoid confusion with a hand-written u it was replaced by "ij". Afterwards the long /i:/ sound it represented became a diphthong /ɛɪ̯/ (although many dialects retain /i/). To write the /i/ sound Dutch mostly uses -ie. The covered version as in "rit" corresponds more to the English sound in "will" or "rid", than to the German sound in "mit" or "Kind". As said above, the distinction 'short'-'long' has little to do with pure length, because the change from open to closed is much more important. There is an exception. In front of -r the long vowel may indeed just be the same vowel held a bit longer: In front of -r there are a few other oddities: Diphthongs. Diphthongs like /ɔɪ̯/ (as in English "toy") or /ɑɪ̯/ (as in English "my") are not used in the standard language. In various dialects they do occur and producing them is often frowned upon. They are considered 'lower class' in many circles. In unstressed syllables like the suffix -lijk the ij represents a schwa. Medeklinkers. Most consonants in Dutch are pronounced more or less the same way as in English but there are a number of notable exceptions. First of all a number of phonemes that English has are simply missing in Dutch. Phonemes are sounds that suffice in marking one word as different from the other. Missing phonemes. Please avoid these sounds when speaking Dutch. Even /ʃ/ : "sh" as in "sh"ip is not really a native Dutch sound, but it occurs in quite a few loans from various sources (Frisian, English, German, French) and as most Dutch people learn English these days they are quite familiar with it. The same thing goes for /ʒ/ as in garage and their affricate versions as in /ʧ/ in church or /ʤ/ George. None of them are native to Dutch. The g, ch and sch problem. The spelling sch- can be rather confusing for people familiar with some German. In German it is used to write the /ʃ/ sound, where English uses sh-. In Dutch the sch- combination also occurs quite frequently but is pronounced rather differently. In most cases it presents a combination of s+ch where the latter is the voiceless velar fricative /x/ as heard in German "Bach" or Scottish "loch". In older versions of the orthography (prior to 1947) the combination -sch represented a simple /s/ sound in final position. The guttural ch at the end had gone mute. (Originally it represented a k- sound as it still does in some dialects and in Frisian). The final -sch spelling is still used for one rather common ending: -isch and also in numerous geographical names as they have never been altered in spelling: In principle Dutch has both a voiceless and a voiced velar fricative and the letter 'g' represents the voiced one and the combination 'ch' the voiceless one. However, the number of words where this creates a phonemic distinction is very small: It depends on the region whether this distinction is actually made in the spoken language. Around Amsterdam it would not be, further south the phoneme 'g' is often pronounced as a voiced palatal fricative, so that the difference becomes more pronounced. Worldwide, the voiced /ɣ/ sound is pretty rare. It only occurs in a few languages like Arabic and Gaelic. As many native speakers do not use it either, it is recommended not to bother about it and use the voiceless /x/ for both, unless your mother tongue happens to have the difference. The Dutch "r". Another, similar, problem for English-speaking learners is the Dutch "r". Essentially there are two, both of which were historically trilled. Alveolar. The first, alveolar /r/, is the "rolling" trilled "r" also used in Spanish, produced with the tip of the tongue against the alveolar ridge. This sound was considered standard only a few decades ago and is still used by quite a few speakers. That includes the author of this book, who grew up at the southern edge of the Randstad, the Drechtsteden region, as well as many people outside the Randstad, including Flanders. Uvular. However, particularly in the Randstad, the rolling alveolar /r/ is gradually replaced by a voiced velar or uvular trill [ʁ], which is also used in many dialects of German; it is also similar to the French "r", but is voiced and articulated somewhat further forward (it is less "throaty"). Both /r/ and /ʁ/ sounds are often not trilled when spoken by native speakers; the alveolar [r] is more often a light tap, while the uvular /ʁ/ can turn into a fricative or approximant. This also presents a considerable challenge for those unaccustomed to the sound when they are confronted with words like "groot" ("big"). The first two sounds tend to blend to one lengthy velar/uvular /xoːt/ or /ɣoːt/, which may cause confusion with words like "rood" (red) and "goot" (gutter). Approximant. A third type of r currently making inroads into Dutch is "Gooise R" (/ɹ/), named after the Gooi area in the Netherlands, where Dutch television is produced (Hilversum) and this speech feature is popularly thought to have originated. This concerns an approximant sound, not unlike American /ɹ/ in words such as "bar", without trilling or friction. Voiceless consonants. The latter word also contains an exception on the rule that t represents /t/. In the ending -tie (corresponding to -tion) it is pronounced as a quick /ʦ/ combination, as in Russian "ц". In contrast to English it is not silent in combinations like kn-: "knie"' /kni/ Devoicing and assimilation. As in German, but unlike English all consonants at the ends of words are devoiced ("..at the ents off worts are devoist..."). You may hear that phenomenon when people speak English with a strong Dutch accent. Assimilation with the previous word often devoices the consonant in initial position as well: The neutral article "het" is often reduced to a prefixed t-sound in the spoken language and occasionally rendered as such in the written language as: 't. Het zaad -&gt; 't zaad. Notice however that both the /z/ and the /d/ reappear in the plural: Contrary to d, the letters v and z are not written in the final position in such cases: Voiced consonants. Apart from the devoicing effects Dutch has the following voiced consonants: Around Amsterdam the tendency to devoice is so strong that /v/ /ɣ/ and /z/ are seldom heard. People use /f/, /x/ and /s/ instead. Liquids, nasals, clusters. Thus, in the Netherlands f, v and w are all pronounced with the upper teeth resting on the lower lips but there is a distinction in voicing and in aspiration (blowing) In "erwt" /ɛrt/ (pea) the w is silent for most speakers, but in initial wr- ("wraak" /vraːk/ or /ʋraːk/) w tends to sound more like /v/. Although the pronunciation of "w" varies, it is "not silent" as in English. For the combination kn- the same holds: the 'k' is clearly pronounced: Many plurals (including the plural forms of the verb) have an ending -en. For many speakers this is pronounced as /-ə/ and the final n is dropped, but this is not true for all speakers. It depends strongly on the region of the speech area you are in. Around Amsterdam it is certainly /-ə/, but in Groningen or in West-Flanders it is a syllabic /-n/ instead. It also depends on how well people are trying to articulate or on the next word. If the latter starts with a vowel the -n is more likely to be pronounced. Even in loans the /ŋg/ tends to be avoided: mango : [ˈmɑnɣo] or [ˈmɑŋɣo] In combination with i it forms a diphthong: ij. Although this is a two letter combination, both letters get capitalized at the beginning of a sentence: ijs -&gt; IJs (ice). /ɛis/. The suffix -je that forms the rather ubiquitous diminutives tends to palatellize the previous consonants or even fuse to a palatal stop all together in rapid speech. It is much less used than in German, e.g. theater: /teˈjatər/ rather than /teˈʔatər/. Clusters of consonants are common in Dutch although perhaps less so than in language like Russian. Usually speakers will pronounce all the consonants in the cluster, but if clusters are consecutive, e.g. in compound words, some elision may occur. E.g. in the compound "vaststellen" typically only one 'st' is actually pronounced. Clusters can occur both initially and finally in a syllable: Common initial clusters: bl-, br-, chl- chr-, dr-, dw-, fl-, fn-, fr-, gl-, gn-, gr-, kl-, kn-, kr-, kw-, pl-, pr-, sch-, schr-, sf-, sj-, sk-, sl-, sm-, sn-, sp-, spl-, spr-, tr-, tw-, vl-, vr-, wr-, zw- Syllable Stress. Dutch is -like English and German but unlike French- a typical stress language. One syllable tends to get all the attention. It is at the same time loud, long and high in pitch and it never has a schwa ə. Instead it has a full vowel or a diphthong. Unstressed syllables tend to be short, low, soft and usually have a schwa, although there are exceptions. Stress is not represented in the spelling unless ambiguity can arise for native speakers. In that case acutes can be added, otherwise orthographic rules demand that they be omitted. For non-native speakers this is a bit of a problem, but often an educated guess can be made which syllable is the stressed one. Because the schwa is written as a "e" in the orthography it is often quite clear where the stress falls in a word: Unfortunately for the non-native speaker the letter e is also used for other purposes. The middle -e represent a full /e/ sound (as the ai in bait), but that is only clear for a native speaker. But with a bit of knowledge of grammar it will be clear that /led/ is the root of a verb (lijden actually) and that ver- is a prefix and -en the suffix. In general, the root of a verb (or noun) will get the stress in Dutch. Of course, there are complicated cases: It's a word which has three e's, each pronounced differently. The first syllable has a full [e] even though it is not stressed. Of course there are cases where stress is not clear even for native speakers; a good example are the separable and non-separable verbs, e.g. the verb doorlopen can be either dóórlopen or doorlópen with a different conjugation and different meaning. In such cases Dutch spelling does allow stress patterns to be written with an acute, but only if otherwise confusion might arise. Some words and names can have rather surprising stress patterns: Capitalization. The rules for capitalization in Dutch are similar to those in English. Capitalization occurs at the beginning of a sentence. "Eigennamen" (proper names, e.g., of persons, organisations, countries etc.) are capitalized, but "soortnamen" (generic names) are not. As mentioned above, when a word beginning with ij has to be capitalized, both letters become capitals, e.g. IJsselmeer. 

Introduction. In this chapter, we start to practice working with XML using XML documents, schemas, and stylesheets. An XML document organizes data and information in a structured, hierarchical format. An XML schema provides standards and rules for the structure of a given XML document. An XML schema also enables data transfer. An XSL (XML stylesheet) allows unique presentations of the material found within an XML document. In the first chapter, "Introduction to XML", you learned what XML is, why it is useful, and how it is used. So, now you want to create your very own XML documents. In this chapter, we will show you the basic components used to create an XML document. This chapter is the foundation for all subsequent chapters--it is a little lengthy, but don't be intimidated. We will take you through the fundamentals of XML documents. This chapter is divided into three parts: As you learned in the previous chapter, the XML Schema and Stylesheet are essentially specialized XML Documents. Within each of these three parts we will examine the layout and components required to create the document. There are links at the end of the XML document, schema, and stylesheet sections that show you how to create the documents using an XML editor. At the bottom of the page there is a link to Exercises for this chapter and a link to the Answers. The first thing you will need before starting to create XML documents is a problem--something you want to solve by using XML to store and share data or information. You need some entity you can collect information about and then access in a variety of formats. So, we created one for you. To develop an XML document and schema, start with a data model depicting the reality of the actual data that is exchanged. Once a high fidelity model has been created, the data model can be readily converted to an XML document and schema. In this chapter, we start with a very simple situation and in successive chapters extend the complexity to teach you more features of XML. Our starting point is a single entity, CITY, which is shown in the following figure. While our focus is on this single entity, to map CITY to an XML schema, we need to have an entity that contains CITY. In this case, we have created TOURGUIDE. Think of a TOURGUIDE as containing many cities, and in this case TOURGUIDE has no attributes nor an identifier. It is just a container for data about cities. Exhibit 1: Data model - Tourguide XML document. An XML document is a file containing XML code and syntax. XML documents have an .xml file extension. We will examine the features &amp; components of the XML document. Below is a sample XML document using our TourGuide model. We will refer to it as we describe the parts of an XML document. Exhibit 2: XML document for city entity  &lt;?xml version="1.0" encoding="UTF-8"?&gt;  &lt;tourGuide xmlns:xsi='http://www.w3.org/2001/XMLSchema-instance'  xsi:noNamespaceSchemaLocation='city.xsd'&gt;  &lt;city&gt;  &lt;cityName&gt;Belmopan&lt;/cityName&gt;  &lt;adminUnit&gt;Cayo&lt;/adminUnit&gt;  &lt;country&gt;Belize&lt;/country&gt;  &lt;population&gt;11100&lt;/population&gt;  &lt;area&gt;5&lt;/area&gt;  &lt;elevation&gt;130&lt;/elevation&gt;  &lt;longitude&gt;88.44&lt;/longitude&gt;  &lt;latitude&gt;17.27&lt;/latitude&gt;  &lt;description&gt;Belmopan is the capital of Belize&lt;/description&gt;  &lt;history&gt;Belmopan was established following the devastation of the  former capital, Belize City, by Hurricane Hattie in 1965. High  ground and open space influenced the choice and ground-breaking  began in 1966. By 1970 most government offices and operations had  already moved to the new location.  &lt;/history&gt;  &lt;/city&gt;  &lt;city&gt;  &lt;cityName&gt;Kuala Lumpur&lt;/cityName&gt;  &lt;adminUnit&gt;Selangor&lt;/adminUnit&gt;  &lt;country&gt;Malaysia&lt;/country&gt;  &lt;population&gt;1448600&lt;/population&gt;  &lt;area&gt;243&lt;/area&gt;  &lt;elevation&gt;111&lt;/elevation&gt;  &lt;longitude&gt;101.71&lt;/longitude&gt;  &lt;latitude&gt;3.16&lt;/latitude&gt;  &lt;description&gt;Kuala Lumpur is the capital of Malaysia and the largest  city in the nation&lt;/description&gt;  &lt;history&gt;The city was founded in 1857 by Chinese tin miners and  preceded Klang. In 1880 the British government transferred their  headquarters from Klang to Kuala Lumpur, and in 1896 it became the  capital of Malaysia.  &lt;/history&gt;  &lt;/city&gt;  &lt;city&gt;  &lt;cityName&gt;Winnipeg&lt;/cityName&gt;  &lt;adminUnit&gt;St. Boniface&lt;/adminUnit&gt;  &lt;country&gt;Canada&lt;/country&gt;  &lt;population&gt;618512&lt;/population&gt;  &lt;area&gt;124&lt;/area&gt;  &lt;elevation&gt;40&lt;/elevation&gt;  &lt;longitude&gt;97.14&lt;/longitude&gt;  &lt;latitude&gt;49.54&lt;/latitude&gt;  &lt;description&gt;Winnipeg has two seasons. Winter and Construction.&lt;/description&gt;  &lt;history&gt;The city was founded by people at the forks (Fort Garry)  trading in pelts with the Hudson Bay Company. Ironically,  The Bay was bought by America.  &lt;/history&gt;  &lt;/city&gt;  &lt;/tourGuide&gt; Prologue (XML declaration). The XML document starts off with the prologue. The prologue informs both a reader and the computer of certain specifications that make the document XML compliant. The first line is the XML declaration (and the only line in this basic XML document). Exhibit 3: XML document - prologue  &lt;?xml version="1.0" encoding="UTF-8"?&gt; xml   =   this is an XML document version="1.0"   =   the XML version (XML 1.0 is the W3C-recommended version) encoding="UTF-8"   =   the character encoding used in the document - UTF 8 corresponds to 8-bit encoded Unicode characters (i.e. the standard way to encode international documents) - Unicode provides a unique number for every character. Another potential attribute of the XML declaration: standalone="yes"   =   the dependency of the document ('yes' indicates that the document does not require another document to complete content) Elements. The majority of what you see in the XML document consists of XML elements. Elements are identified by their tags that open with &lt; or &lt;/ and close with &gt; or /&gt;. The start tag looks like this: &lt;element attribute="value"&gt;, with a left angle bracket (&lt;) followed by the element type name, optional attributes, and finally a right angle bracket (&gt;). The end tag looks like this: &lt;/element&gt;, similar to the start tag, but with a slash (/) between the left angle bracket and the element type name, and no attributes. When there's nothing between a start tag and an end tag, XML allows you to combine them into an empty element tag, which can include everything a start tag can: &lt;img src="Belize.gif" /&gt;. This one tag must be closed with a slash and right angle bracket (/&gt;), so that it can be distinguished from a start tag. The XML document is designed around a major theme, an umbrella concept covering all other items and subjects; this theme is analyzed to determine its component parts, creating categories and subcategories. The major theme and its component parts are described by elements. In our sample XML document, 'tourGuide' is the major theme; 'city' is a category; 'population' is a subcategory of 'city'; and the hierarchy may be carried even further: 'males' and 'females' could be subcategories of 'population'. Elements follow several rules of syntax that will be described in the Rules to Follow section. We left out the attributes within the &lt;tourGuide&gt; start tag — that part will be explained in the XML Schema section. Exhibit 4: Elements of the city entity XML document  &lt;tourGuide&gt;  &lt;city&gt;  &lt;cityName&gt;Belmopan&lt;/cityName&gt;  &lt;adminUnit&gt;Cayo&lt;/adminUnit&gt;  &lt;country&gt;Belize&lt;/country&gt;  &lt;population&gt;11100&lt;/population&gt;  &lt;area&gt;5&lt;/area&gt;  &lt;elevation&gt;130&lt;/elevation&gt;  &lt;longitude&gt;88.44&lt;/longitude&gt;  &lt;latitude&gt;17.27&lt;/latitude&gt;  &lt;description&gt;Belmopan is the capital of Belize&lt;/description&gt;  &lt;history&gt;Belmopan was established following the devastation of the  former capital, Belize City, by Hurricane Hattie in 1965. High  ground and open space influenced the choice and ground-breaking  began in 1966. By 1970 most government offices and operations had  already moved to the new location.  &lt;/history&gt;  &lt;/city&gt;  &lt;/tourGuide&gt; Element hierarchy. &lt;br&gt; Attributes. Attributes aid in modifying the content of a given element by providing additional or required information. They are contained within the element's opening tag. In our sample XML document code we could have taken advantage of attributes to specify the unit of measure used to determine the area and the elevation (it could be feet, yards, meters, kilometers, etc.); in this case, we could have called the attribute 'measureUnit' and defined it within the opening tag of 'area' and 'elevation'.  &lt;adminUnit class="state"&gt;Cayo&lt;/adminUnit&gt;  &lt;adminUnit class="region"&gt;Selangor&lt;/adminUnit&gt; The above attribute example can also be written as: 1. using child elements  &lt;adminUnit&gt;  &lt;class&gt;state&lt;/class&gt;  &lt;name&gt;Cayo&lt;/name&gt;  &lt;/adminUnit&gt;  &lt;adminUnit&gt;  &lt;class&gt;region&lt;/class&gt;  &lt;name&gt;Selangor&lt;/name&gt;  &lt;/adminUnit&gt; 2. using an empty element  &lt;adminUnit class="state" name="Cayo" /&gt;  &lt;adminUnit class="region" name="Selangor" /&gt; Attributes can be used to: Attributes can, however, be a bit more difficult to manipulate and they have some constraints. Consider using a child element if you need more freedom. Rules to follow. These rules are designed to aid the computer reading your XML document. (e.g. &lt;element&gt;data stuff&lt;/element&gt;). =&gt; the parent element's opening and closing tags must contain all of its child elements' tags; in this way, you close first the tag that was opened last:  &lt;parentElement&gt;        &lt;childElement1&gt;data&lt;/childElement1&gt;        &lt;childElement2&gt;                &lt;subChildElementA&gt;data&lt;/subChildElementA&gt;                &lt;subChildElementB&gt;data&lt;/subChildElementB&gt;        &lt;/childElement2&gt;        &lt;childElement3&gt;data&lt;/childElement3&gt;  &lt;/parentElement&gt; XML Element Naming Convention. Any name can be used but the idea is to make names meaningful to those who might read the document. XML documents often have a corresponding database. The database will contain fields which correspond to elements in the XML document. A good practice is to use the naming rules of your database for the elements in the XML documents. DTD (Document Type Definition) Validation - Simple Example. Simple Internal DTD.  &lt;?xml version="1.0"?&gt;  &lt;!DOCTYPE cdCollection [  &lt;!ELEMENT cdCollection (cd)&gt;  &lt;!ELEMENT cd (title, artist, year)&gt;  &lt;!ELEMENT title (#PCDATA)&gt;  &lt;!ELEMENT artist (#PCDATA)&gt;  &lt;!ELEMENT year (#PCDATA)&gt;  &lt;cdCollection&gt;  &lt;cd&gt;  &lt;title&gt;Dark Side of the Moon&lt;/title&gt;  &lt;artist&gt;Pink Floyd&lt;/artist&gt;  &lt;year&gt;1973&lt;/year&gt;  &lt;/cd&gt;  &lt;/cdCollection&gt; Every element that will be used MUST be included in the DTD. Don’t forget to include the root element, even though you have already specified it at the beginning of the DTD. You must specify it again, in an &lt;!ELEMENT&gt; tag. &lt;!ELEMENT cdCollection (cd)&gt; The root element, &lt;cdCollection&gt;, contains all the other elements of the document, but only one direct child element: &lt;cd&gt;. Therefore, you need to specify the child element (only direct child elements need to be specified) in the parentheses. &lt;!ELEMENT cd (title, artist, year)&gt; With this line, we define the &lt;cd&gt; element. Note that this element contains the child elements &lt;title&gt;, &lt;artist&gt;, and &lt;year&gt;. These are spelled out in a particular order. This order must be followed when creating the XML document. If you change the order of the elements (with this particular DTD), the document won’t validate. &lt;!ELEMENT title (#PCDATA)&gt; The remaining three tags, &lt;title&gt;, &lt;artist&gt;, and &lt;year&gt; don’t actually contain other tags. They do however contain some text that needs to be parsed. You may remember from an earlier lecture that this data is called Parsed Character Data, or #PCDATA. Therefore, #PCDATA is specified in the parentheses. So this simple DTD outlines exactly what you see here in the XML file. Nothing can be added or taken away, as long as we stick to this DTD. The only thing you can change is the #PCDATA text part between the tags. Adding complexity. There may be times when you will want to put more than just character data, or more than just child elements into a particular element. This is referred to as mixed content. For example, let’s say you want to be able to put character data OR a child element, such as the tag into a &lt;description&gt; element:  &lt;!ELEMENT description (#PCDATA | b | i )*&gt; This particular arrangement allows us to use PCDATA, the tag, or the tag all at once. One particular caveat though, is that if you are going to mix PCDATA and other elements, the grouping must be followed by the asterisk (*) suffix. This declaration allows us to now add the following to the XML document (after defining the individual elements of course)  &lt;cd&gt;  &lt;title&gt;Love. Angel. Music. Baby&lt;/title&gt;  &lt;artist&gt;Gwen Stefani&lt;/artist&gt;  &lt;year&gt;2004&lt;/year&gt;  &lt;genre&gt;pop&lt;/genre&gt;  &lt;description&gt;  This is a great album from former  &lt;nowiki&gt;&lt;i&gt;No Doubt&lt;/i&gt; singer &lt;b&gt;Gwen Stephani&lt;/b&gt;.&lt;/nowiki&gt;  &lt;/description&gt;  &lt;/cd&gt; With attributes this is done a little differently than with elements. Please see following example:  &lt;cd remaster_date=”1992”&gt;  &lt;title&gt;Dark Side of the Moon&lt;/title&gt;  &lt;artist&gt;Pink Floyd&lt;/artist&gt;  &lt;year&gt;1973&lt;/year&gt;  &lt;/cd&gt; In order for this to validate, it must be specified in the DTD. Attribute content models are specified with:  &lt;!ATTLIST element_name attribute_name attribute_type default_value&gt; Let’s use this to validate our CD example:  &lt;!ATTLIST cd remaster_date CDATA #IMPLIED&gt; Choices.  &lt;ATTLIST person gender (male|female) “male”&gt; Grouping Attributes for an Element. If a particular element is to have many different attributes, group them together like so:  &lt;!ATTLIST car horn CDATA #REQUIRED  seats CDATA #REQUIRED  steeringwheel CDATA #REQUIRED  price CDATA #IMPLIED&gt; Adding STATIC validation, for items that must have a certain value.  &lt;!ATTLIST classList classNumber CDATA #IMPLIED  building (UWINNIPEG_DCE|UWINNIPEG_MAIN) "UWINNIPEG_MAIN"  originalDeveloper CDATA #FIXED "Khal Shariff"&gt; Suffixes. So what happens with our last example with the CD collection, when we want to add more CDs? With the current DTD, we cannot add any more CDs without getting an error. Try it and see. When you specify a child element (or elements) the way we did, only one of each child element can be used. Not very suitable for a CD collection is it? We can use something called suffixes to add functionality to the &lt;!ELEMENT&gt; tag. Suffixes are added to the end of the specified child element(s). There are 3 main suffixes that can be used: Validating for multiple children with a DTD. So in the case of our CD collection XML file, we can add more CDs to the list by adding a + suffix:  &lt;!ELEMENT cd_collection(cd+)&gt; Using more internal formatting tags. Bold tags, B's for example are also defined in the DTD as elements, that are optional like thus:  &lt;ELEMENT notes (#PCDATA | b | i)*&gt;  &lt;!ELEMENT b (#PCDATA)*&gt;  &lt;!ELEMENT i (#PCDATA)*&gt; _______________  &lt;classList classNumber="303" building="UWINNIPEG_DCE" originalDeveloper="Khal Shariff"&gt;  &lt;student&gt;  &lt;firstName&gt;Kenneth  &lt;/firstName&gt;  &lt;lastName&gt;Branaugh  &lt;/lastName&gt;  &lt;studentNumber&gt;  &lt;/studentNumber&gt;  &lt;notes&gt;Excellent , Kenneth is doing well.  &lt;/notes&gt;  etc ONLINE Validator. http://www.stg.brown.edu/service/xmlvalid/ Well-formed and valid XML. Well-formed XML  -  An XML document that correctly abides by the rules of XML syntax. Valid XML  -  An XML document that adheres to the rules of an XML schema (which we will discuss shortly). To be valid an XML document must first be well-formed. A Valid XML Document must be Well-formed. But, a Well-formed XML Document might not be valid - in other words, a well-formed XML document, that meets the criteria for XML syntax, might not meet the criteria for the XML schema, and will therefore be invalid. For example, think of the situation where your XML document contains the following (for this schema):  &lt;city&gt;  &lt;cityName&gt;Boston&lt;/cityName&gt;  &lt;country&gt;United States&lt;/country&gt;  &lt;adminUnit&gt;Massachusetts&lt;/adminUnit&gt;  :  &lt;/city&gt; Notice that the elements do not appear in the correct sequence according to the schema (cityName, adminUnit, country). The XML document can be validated (using validation software) against its declared schema – the validation software would then catch the out of sequence error. Using an XML Editor. Check chapter ../XML Editor/ for instructions on how to start an XML editor. Once you have followed the steps to get started you can copy the code in the sample XML document and paste it into the XML editor. Then check your results. Is the XML document well-formed? Is the XML document valid? (you will need to have copied and pasted the schema in order to validate - we will look at schemas next) XML schema. An XML schema is an XML document. XML schemas have an .xsd file extension. An XML schema is used to govern the structure and content of an XML document by providing a template for XML documents to follow in order to be valid. It is a guide for how to structure your XML document as well as indicating your XML document's components (elements and attributes - and their relationships). An XML editor will examine an XML document to ensure that it conforms to the specifications of the XML schema it is written against - to ensure it is valid. XML schemas engender confidence in data transfer. With schemas, the receiver of data can feel confident that the data conforms to expectations. The sender and the receiver have a mutual understanding of what the data represent. Because an XML schema is an XML document, you use the same language - standard XML markup syntax - with elements and attributes specific to schemas. A schema defines: For more detailed information on XML schemas and reference lists of: Common XML Schema Primitive Data Types, Summary of XML Schema Elements, Schema Restrictions and Facets for data types, and Instance Document Attributes, click on this wikibook link =&gt; http://en.wikibooks.org/wiki/XML_Schema Schema reference. This is the part of the XML Document that references an XML Schema: Exhibit 5: XML document's schema reference  &lt;tourGuide  xmlns:xsi='http://www.w3.org/2001/XMLSchema-instance'  xsi:noNamespaceSchemaLocation='city.xsd'&gt; This is the part we left out when we described the root element in the basic XML document from the previous section. The additional attributes of the root element &lt;tourGuide&gt; reference the XML schema (it is the schemaLocation attribute). xmlns:xsi='http://www.w3.org/2001/XMLSchema-instance'  -  references the W3C Schema-instance namespace xsi:noNamespaceSchemaLocation='city.xsd'  -  references the XML schema document (city.xsd) Schema document. Below is a sample XML schema using our TourGuide model. We will refer to it as we describe the parts of an XML schema. Exhibit 6: XML schema document for city entity  &lt;?xml version="1.0" encoding="UTF-8"?&gt;  &lt;xsd:schema xmlns:xsd="http://www.w3.org/2001/XMLSchema"  elementFormDefault="unqualified"&gt;  &lt;xsd:element name="tourGuide"&gt;  &lt;xsd:complexType&gt;  &lt;xsd:sequence&gt;  &lt;xsd:element name="city" type="cityDetails" minOccurs = "1" maxOccurs="unbounded" /&gt;  &lt;/xsd:sequence&gt;  &lt;/xsd:complexType&gt;  &lt;/xsd:element&gt;  &lt;xsd:complexType name="cityDetails"&gt;  &lt;xsd:sequence&gt;  &lt;xsd:element name="cityName" type="xsd:string"/&gt;  &lt;xsd:element name="adminUnit" type="xsd:string"/&gt;  &lt;xsd:element name="country" type="xsd:string"/&gt;  &lt;xsd:element name="population" type="xsd:integer"/&gt;  &lt;xsd:element name="area" type="xsd:integer"/&gt;  &lt;xsd:element name="elevation" type="xsd:integer"/&gt;  &lt;xsd:element name="longitude" type="xsd:decimal"/&gt;  &lt;xsd:element name="latitude" type="xsd:decimal"/&gt;  &lt;xsd:element name="description" type="xsd:string"/&gt;  &lt;xsd:element name="history" type="xsd:string"/&gt;  &lt;/xsd:sequence&gt;  &lt;/xsd:complexType&gt;  &lt;/xsd:schema&gt;  Note: Latitude and Longitude are decimal data types.  The conversion is from the usual form (e.g., 50º 17' 35")  to a decimal by using the formula degrees+min/60+secs/3600. Prolog. Remember that the XML schema is essentially an XML document and therefore must begin with the prolog, which in the case of a schema includes: &lt;br&gt; The XML declaration:  &lt;?xml version="1.0" encoding="UTF-8"?&gt; The schema element declaration: &lt;xsd:schema xmlns:xsd="http://www.w3.org/2001/XMLSchema" elementFormDefault="unqualified"&gt; The schema element is similar to a root element - it contains all other elements in the schema. Attributes of the schema element include: xmlns  -  XML NameSpace - the URL for the site that describes the XML elements and data types used in the schema. You can find more about namespaces here =&gt; ../Namespace/. xmlns:xsd  -  All the elements and attributes with the 'xsd' prefix adhere to the vocabulary designated in the given namespace. elementFormDefault  -  elements from the target namespace are either required or not required to be qualified with the namespace prefix. This is mostly useful when more than one namespace is referenced. In this case, 'elementFormDefault' must be "qualified", because you must indicate which namespace you are using for each element. If you are referencing only one namespace, then 'elementFormDefault' can be "unqualified". Perhaps, using "qualified" as the default is most prudent, this way you do not accidentally forget to indicate which namespace you are referencing. Element declarations. Define the elements in the schema. Include: Basic element declaration format: &lt;xsd:element name="name" type="type"&gt; Simple type. declares elements that: example: &lt;xsd:element name="cityName" type="xsd:string" /&gt; Default Value If an element is not assigned a value then the default value is assigned. example: &lt;xsd:element name="description" type="xsd:string" default="really cool place to visit!" /&gt; Fixed Value An attribute that is defined as fixed must be empty or contained the specified fixed value. No other values are allowed. example: &lt;xsd:element name="description" type="xsd:string" fixed="you must visit this place - it is awesome!" /&gt; Complex type. declares elements that: examples: 1. The root element 'tourGuide' contains a child element 'city'. This is shown here: Nameless complex type  &lt;xsd:element name="tourGuide"&gt;  &lt;xsd:complexType&gt;  &lt;xsd:sequence&gt;  &lt;xsd:element name="city" type="cityDetails" minOccurs = "1" maxOccurs="unbounded" /&gt;  &lt;/xsd:sequence&gt;  &lt;/xsd:complexType&gt;  &lt;/xsd:element&gt; Occurrence Indicators: 2. The parent element 'city' contains many child elements: 'cityName', 'adminUnit', 'country', 'population', etc. Why does this complex element set not start with the line: &lt;xsd:element name="city" type="cityDetails"&gt;? The element 'city' was already defined above within the complex element 'tourGuide' and it was given the type, 'cityDetails'. This data type, 'cityDetails', is utilized here in identifying the sequence of child elements for the parent element 'city'. Named Complex Type - and therefore can be reused in other parts of the schema  &lt;xsd:complexType name="cityDetails"&gt;  &lt;xsd:sequence&gt;  &lt;xsd:element name="cityName" type="xsd:string"/&gt;  &lt;xsd:element name="adminUnit" type="xsd:string"/&gt;  &lt;xsd:element name="country" type="xsd:string"/&gt;  &lt;xsd:element name="population" type="xsd:integer"/&gt;  &lt;xsd:element name="area" type="xsd:integer"/&gt;  &lt;xsd:element name="elevation" type="xsd:integer"/&gt;  &lt;xsd:element name="longitude" type="xsd:decimal"/&gt;  &lt;xsd:element name="latitude" type="xsd:decimal"/&gt;  &lt;xsd:element name="description" type="xsd:string"/&gt;  &lt;xsd:element name="history" type="xsd:string"/&gt;  &lt;/xsd:sequence&gt;  &lt;/xsd:complexType&gt; The &lt;xsd:sequence&gt; tag indicates that the child elements must appear in the order, the sequence, specified here. Compare the sample XML Schema and the sample XML Document - try to observe patterns in the code and how the XML Schema sets up the XML Document. 3. Elements that have attributes are also designated as complex type. a. this XML Document line: &lt;adminUnit class="state" name="Cayo" /&gt; would be defined in the XML Schema as:  &lt;xsd:element name="adminUnit"&gt;  &lt;xsd:complexType&gt;  &lt;xsd:attribute name="class" type="xsd:string" /&gt;  &lt;xsd:attribute name="name" type="xsd:string" /&gt;  &lt;/xsd:complexType&gt;  &lt;/xsd:element&gt; b. this XML Document line: &lt;adminUnit class="state"&gt;Cayo&lt;/adminUnit&gt; would be defined in the XML Schema as:  &lt;xsd:element name="adminUnit"&gt;  &lt;xsd:complexType&gt;  &lt;xsd:simpleContent&gt;  &lt;xsd:extension base="xsd:string"&gt;  &lt;xsd:attribute name="class" type="xsd:string" /&gt;  &lt;/xsd:extension&gt;  &lt;/xsd:simpleContent&gt;  &lt;/xsd:complexType&gt;  &lt;/xsd:element&gt; Attribute declarations. Attribute declarations are used in complex type definitions. We saw some attribute declarations in the third example of the Complex Type Element. &lt;xsd:attribute name="class" type="xsd:string" /&gt; Data type declarations. These are contained within element and attribute declarations as: type=" " . Common XML Schema Data Types XML schema has a lot of built-in data types. The most common types are: For an entire list of built-in simple data types see http://www.w3.org/TR/xmlschema-2/#built-in-datatypes Using an XML Editor =&gt; ../XML Editor/ This link will take you to instructions on how to start an XML editor. Once you have followed the steps to get started you can copy the code in the sample XML schema document and paste it into the XML editor. Then check your results. Is the XML schema well-formed? Is the XML schema valid? XML stylesheet (XSL). An XML Stylesheet is an XML Document. XML Stylesheets have an .xsl file extension. The eXtensible Stylesheet Language (XSL) provides a means to transform and format the contents of an XML document for display. Since an XML document does not contain tags a browser understands, such as HTML tags, browsers cannot present the data without a stylesheet that contains the presentation information. By separating the data and the presentation logic, XSL allows people to view the data according to their different needs and preferences. The XSL Transformation Language (XSLT) is used to transform an XML document from one form to another, such as creating an HTML document to be viewed in a browser. An XSLT stylesheet consists of a set of formatting instructions that dictate how the contents of an XML document will be displayed in a browser, with much the same effect as Cascading Stylesheets (CSS) do for HTML. Multiple views of the same data can be created using different stylesheets. The output of a stylesheet is not restricted to a browser. During the transformation process, XSLT analyzes the XML document and converts it into a node tree – a hierarchical representation of the entire XML document. Each node represents a piece of the XML document, such as an element, attribute or some text content. The XSL stylesheet contains predefined “templates” that contain instructions on what to do with the nodes. XSLT will use the match attribute to relate XML element nodes to the templates, and transform them into the resulting document. Exhibit 7: XML stylesheet document for city entity &lt;?xml version="1.0" encoding="UTF-8"?&gt; &lt;xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"&gt;  &lt;xsl:output method="html"/&gt;  &lt;xsl:template match="/"&gt;  &lt;html&gt;  &lt;head&gt;  &lt;title&gt;Tour Guide&lt;/title&gt;  &lt;/head&gt;  &lt;body&gt;  &lt;h2&gt;Cities&lt;/h2&gt;  &lt;xsl:apply-templates select="tourGuide"/&gt;  &lt;/body&gt;  &lt;/html&gt;  &lt;/xsl:template&gt;  &lt;xsl:template match="tourGuide"&gt;  &lt;xsl:for-each select="city"&gt;  &lt;br/&gt;&lt;xsl:value-of select="continentName"/&gt;&lt;br/&gt;  &lt;xsl:value-of select="cityName"/&gt;&lt;br/&gt;  &lt;xsl:text&gt;Population: &lt;/xsl:text&gt;  &lt;xsl:value-of select='format-number(population, "##,###,###")'/&gt;&lt;br/&gt;  &lt;xsl:value-of select="country"/&gt;  &lt;br/&gt;  &lt;/xsl:for-each&gt;  &lt;/xsl:template&gt; &lt;/xsl:stylesheet&gt; The output of the city.xsl stylesheet in Table 2-3 will look like the following: You will notice that the stylesheet consists of HTML to inform the media tool (a web browser) of the presentation design. If you do not already know HTML this may seem a little confusing. Online resources such as the W3Schools tutorials can help with the basic understanding you will need =&gt;(http://www.w3schools.com/html/default.asp). Incorporated within the HTML is the XML that supplies the data, the information, contained within our XML document. The XML of the stylesheet indicates what information will be displayed and how. So, the HTML constructs a display and the XML plugs in values within that display. XSL is the tool that transforms the information into presentational form, but at the same time keeps the meaning of the data. Prolog.  &lt;?xml version="1.0" encoding="UTF-8"?&gt;  &lt;xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"&gt;  &lt;xsl:output method="html"/&gt; The XML declaration  &lt;?xml version="1.0" encoding="UTF-8"?&gt; The stylesheet &amp; namespace declarations  &lt;xsl:stylesheet version="1.0" xmlns:xsl="http://www.w3.org/1999/XSL/Transform"&gt; The output document format  &lt;xsl:output method="html"/&gt; This element designates the format of the output document and must be a child element of &lt;xsl:stylesheet&gt; Templates. The &lt;xsl:template&gt; element is used to create templates that describe how to display elements and their content. Above, in the XSL introduction, we mentioned that XSL breaks up the XML document into nodes and works on individual nodes. This is done with templates. Each template within an XSL describes a single node. To identify which node a given template is describing, use the 'match' attribute. The value given to the 'match' attribute is called a pattern. Remember: (node tree – a hierarchical representation of the entire XML document. Each node represents a piece of the XML document, such as an element, attribute or some text content). Wherever there is branching in the node tree, there is a node. &lt;xsl:template&gt; defines the start of a template and contains rules to apply when a specified node is matched. the match attribute  &lt;xsl:template match="/"&gt; This template match attribute associates the XML document root (/), the whole branch of the XML source document, with the HTML document root. Contained within this template element is the typical HTML markup found at the beginning of any HTML document. This HTML is written to the output. The XSL looks for the root match and then outputs the HTML, which the browser understands.  &lt;xsl:template match="tourGuide"&gt; This template match attribute associates the element 'tourGuide' with the display rules described within this element. Elements. Elements specific to XSL: For more XSL elements =&gt; http://www.w3schools.com/xsl/xsl_w3celementref.asp . Language-Specific Validation and Transformation Methods. PHP Methods of XML Dom Validation. Using the DOM DocumentObjectModel to validate XML and with a DTD DocumentTypeDeclaration and the PHP language on a server and more http://wiki.cc/php/Dom_validation Browser Methods. Place this line of code in your .xml document after the XML declaration (prologue).  &lt;?xml-stylesheet type="text/xsl" href="tourGuide.xsl"?&gt; PHP XML Production.  &lt;?php  $xmlData = "";  mysql_connect('localhost','root',")  or die('Failed to connect to the DBMS');  // make connection to database  mysql_select_db('issd')  or die('Failed to open the requested database');  $result = mysql_query('SELECT * from students') or die('Query to like get the records failed');  if (mysql_num_rows($result)&lt;1){  die (");  $xmlString = "&lt;classlist&gt;\n";  $xmlString .= "\t&lt;student&gt;";  while ($row = mysql_fetch_array($result)) {  $xmlString .= "  \t&lt;firstName&gt;  ".$row['firstName']."  &lt;/firstName&gt;\n  \t&lt;lastName&gt;  ".$row['lastName']."  \t&lt;/lastName&gt;\n";  $xmlString .= "&lt;/student&gt;\n";  $xmlString .= "&lt;/classlist&gt;";  echo $xmlString;  $myFile = "classList.xml"; //any file  $fh = fopen($myFile, 'w') or die("can't open file"); //create filehandler  fwrite($fh, $xmlString); //write the data into the file  fclose($fh); //ALL DONE! PHP Methods of XSLT Transformation. This one is good for PHP5 and wampserver (latest). Please ensure that *xsl* is NOT commented out in the php.ini file.  &lt;?php  // Load the XML source  $xml = new DOMDocument;  $xml-&gt;load('tourguide.xml');  $xsl = new DOMDocument;  $xsl-&gt;load('tourguide.xsl');  // Configure the transformer  $proc = new XSLTProcessor;  $proc-&gt;importStyleSheet($xsl); // attach the xsl rules  echo $proc-&gt;transformToXML($xml); Example 1, Using within PHP itself (use phpInfo() function to check XSLT extension; enable if needed) This example might produce XHTML. Please note it could produce anything defined by the XSL.  &lt;?php  $xhtmlOutput = xslt_create();  $args = array();  $params = array('foo' =&gt; 'bar');  $theResult = xslt_process(  $xhtmlOutput,  'theContentSource.xml',  'theTransformationSource.xsl',  null,  $args,  $params  xslt_free($xhtmlOutput); // free that memory  // echo theResult or save it to a file or continue processing (perhaps instructions) Example 2:  &lt;?php  if (PHP_VERSION &gt;= 5) {  // Emulate the old xslt library functions  function xslt_create() {  return new XsltProcessor();  function xslt_process($xsltproc,  $xml_arg,  $xsl_arg,  $xslcontainer = null,  $args = null,  $params = null) {  // Start with preparing the arguments  $xml_arg = str_replace('arg:', ", $xml_arg);  $xsl_arg = str_replace('arg:', ", $xsl_arg);  // Create instances of the DomDocument class  $xml = new DomDocument;  $xsl = new DomDocument;  // Load the xml document and the xsl template  $xml-&gt;loadXML($args[$xml_arg]);  $xsl-&gt;loadXML($args[$xsl_arg]);  // Load the xsl template  $xsltproc-&gt;importStyleSheet($xsl);  // Set parameters when defined  if ($params) {  foreach ($params as $param =&gt; $value) {  $xsltproc-&gt;setParameter("", $param, $value);  // Start the transformation  $processed = $xsltproc-&gt;transformToXML($xml);  // Put the result in a file when specified  if ($xslcontainer) {  return @file_put_contents($xslcontainer, $processed);  } else {  return $processed;  function xslt_free($xsltproc) {  unset($xsltproc);  $arguments = array(  '/_xml' =&gt; file_get_contents("xml_files/201945.xml"),  '/_xsl' =&gt; file_get_contents("xml_files/convertToSql_new2.xsl")  $xsltproc = xslt_create();  $html = xslt_process(  $xsltproc,  'arg:/_xml',  'arg:/_xsl',  null,  $arguments  xslt_free($xsltproc);  print $html; PHP file writing code.  $myFile = "testFile.xml"; //any file  $fh = fopen($myFile, 'w') or die("can't open file"); //create filehandler  $stringData = "&lt;foo&gt;\n\t&lt;bar&gt;\n\thello\n"; // get a string ready to write  fwrite($fh, $stringData); //write the data into the file  $stringData2 = "\t&lt;/bar&gt;\n&lt;/foo&gt;";  fwrite($fh, $stringData2); //write more data into the file  fclose($fh); //ALL DONE! XML Colors. For use in your stylesheet: these colors can be used for both background and font http://www.w3schools.com/html/html_colors.asp http://www.w3schools.com/html/html_colorsfull.asp http://www.w3schools.com/html/html_colornames.asp Using an XML Editor =&gt; ../XML Editor/ This link will take you to instructions on how to start an XML editor. Once you have followed the steps to get started you can copy the code in the sample XML stylesheet document and paste it into the XML editor. Then check your results. Is the XML stylesheet well-formed? /Definitions/.  XML  SGML  Dan Connelly  RSS  XML Declaration  parent  child  sibling  element  attribute  *Well-formed XML  PCDATA /Exercises/. Exercise 1. a)Using "tourguide" above as a good example, create an XML document whose root is "classlist" . This CLASSLIST is created from a starting point of single entity, STUDENT. Any number of students contain elements: firstname, lastname, emailaddress. 

Table of Contents. Chapters.  __NOEDITSECTION__ 

« Contents Page Introduction to the Invertebrate Zoology Study Guide. This "Study Guide to Invertebrate Zoology" is a textbook at Wikibooks shelved at the section and intended to establish a course of study in the subject of Invertebrate Zoology, mostly utilizing articles found in Wikipedia, with links to other relevant web sites as appropriate. In some cases, portions of the text from "Wikipedia" articles have been used to materially develop introductory text within the Guide. For the new user, it need be pointed out that "Wikipedia" differs from a standard encyclopedia in two important respects: 1) it is a hypertext document, and 2) it is open and editable, and therefore constantly changing. For the student following this or any guide through "Wikipedia" to cover a specific subject, it is recommended that each article (page) be read first in its entirety, before any hyperlinks are followed to other topics or explanations. It is too easy, otherwise, to simply become lost in a maze of links, and miss the main thrust of an article presented as an assignment from the Guide. Because "Wikipedia" is constantly changing (and, it is hoped, improving) the quality of each article encountered will be variable. Some articles are well written and go to adequate depth, whereas others, lacking a proponent, are shallow and incomplete. This short-coming should diminish with time, but can be a problem. One clear advantage to using this Guide linked to a hypertext like "Wikipedia" is the "circular redundancy with serendipity" factor that arises when an article is read and its hyperlinks followed; this factor can be a powerful learning tool. The persistent reader is subjected to a fairly high degree of repetitive reading, often presenting slightly different perspectives on the same general topic, with the result that learning comes from redundancy. At the same time, some hyperlinks lead down less relevant paths, bringing new and unanticipated knowledge. 

This study guide to invertebrate zoology is a textbook at "Wikibooks" filed under and intended to establish a course of study in the subject of Invertebrate Zoology, mostly utilizing articles found in "Wikipedia", with links to other relevant web sites as appropriate. Contributors are encouraged to spend time developing the relevant articles at "Wikipedia", maintaining this textbook as a "guide" to those articles. At some future time, it may be desirable to simply import whole articles over to build a Zoology textbook rather than a "guide book", however the present approach seems to be a good one for getting a useful text at "Wikibooks" based upon the initial author's experience with "Botany". Another Wikibook, "Ecology", follows a similar approach, but requires more text per page, because of less congruence between chapter material and articles at "Wikipedia". The contributors to this book included: 

../Verbs/ | ../Verbs List/ Vosotros. = Los tiempos de los verbos: por simples y compuestos = Siete tiempos compuestos. = Los tiempos de los verbos: por modos = Gramatically, there are different ways of classifying Spanish verb forms (as well as the verbs of any language). The one most people understand is "tense", but another very important classification is the "mood." A "mood" is differentiated from a "tense" in that it does not express time per se (which is what "tense" really means). Unfortunately, "mood" is not really a very good name either because moods do not express moods of the speaker (happy, sad, etc.) The mood is used to classify how the verb is used grammatically in the sentence; that is, how the verb usage works together with the other parts of the sentence. Spanish has the following verb moods: In many cases, along with the name of the mood also goes a tense. A given form of a verb normally has both mood and tense, although when the mood is indicative usually the name of the mood is not stated. This is one thing that makes the concept of moods hard to understand. When we talk about the "present tense" of a verb, we really should say "indicative mood, present tense". This would make it clear that when we say "subjunctive mood, present tense" we're not dealing with anything particularly strange, only a little different. El modo imperativo. Spanish | ../Verbs/ | ../Verbs List/ 

../Verbs/ The imperfect tense communicates: Verbs are conjugated to the imperfect by taking off the last two letters of the infinitive and replacing them with an ending based on -ía or -aba. Conjugating regular verbs in the imperfect. This is probably the easiest verb tense to conjugate. -AR verbs take on the -aba endings. Here is an example: Bailar (to dance)  bailaba: I danced&lt;br&gt;  bailabas: you danced&lt;br&gt;  bailaba: he, she, it, you (formal) danced&lt;br&gt;  bailábamos: we danced &lt;br&gt;  bailabais: you danced &lt;br&gt;  bailaban: they, you (plural, formal) danced &lt;br&gt; -IR and -ER verbs take on -ía endings. Examples: Venir (to come)  venía: I came &lt;br&gt;  venías: you came &lt;br&gt;  venía: he, she, it, you came &lt;br&gt;  veníamos: we came &lt;br&gt;  veníais: you came&lt;br&gt;  venían: they, you (formal) came &lt;br&gt; Comer (to eat) comía: I ate &lt;br&gt; comías: you ate &lt;br&gt; comía: he, she, it, you ate &lt;br&gt; comíamos: we ate &lt;br&gt; comíais: you ate&lt;br&gt; comían: they, you (formal) ate &lt;br&gt; Irregular verbs. There are only three. Ser (to be) era: I was &lt;br&gt; eras: you were &lt;br&gt; era: he, she, it was, you were &lt;br&gt; éramos: we were &lt;br&gt; erais: you were &lt;br&gt; eran: they were &lt;br&gt; Notice the accent in the "nosotros" form of "ser" and "ir." Ir (to go) iba: I used to go&lt;br&gt; ibas: you used to go &lt;br&gt; iba: he, she, it, you used to go &lt;br&gt; íbamos: we used to go &lt;br&gt; ibais: you used to go &lt;br&gt; iban: they used to go &lt;br&gt; Ver (to see) veía: I used to see&lt;br&gt; veías: you used to see &lt;br&gt; veía: he, she, it, you used to see &lt;br&gt; veíamos: we used to see&lt;br&gt; veíais: you used to see&lt;br&gt; veían: they used to see &lt;br&gt; That's it! Those are all the irregulars. 

../Verbs/ The preterite verb form describes past actions that have begun or completed. The preterite sets something in the past as an event, with a beginning and an end. Spanish has another past tense form, the imperfect, that is used for past actions that are continuing, or for which its beginning or its end is not important. Conjugate regular verbs in the preterite. -AR verbs Eliminate the "-ar" and replace it with one of the following endings: -é &lt;br&gt; -aste &lt;br&gt; -ó &lt;br&gt; -amos &lt;br&gt; -asteis&lt;br&gt; -aron &lt;br&gt; Patinar (to skate) becomes: patiné: I skated &lt;br&gt; patinaste: you skated &lt;br&gt; patinó: he skated &lt;br&gt; patinamos: we skated &lt;br&gt; patinasteis: you skated&lt;br&gt; patinaron: they skated &lt;br&gt; Notice that the "nosotros" form conjugates the same as it does in the present tense. This can cause confusion so look for context tools to make sure you have the right tense. -ER and -IR verbs Eliminate the "-er" or "-ir" and replace it with one of the following endings: -í &lt;br&gt; -iste &lt;br&gt; -ió &lt;br&gt; -imos &lt;br&gt; -isteis&lt;br&gt; -ieron &lt;br&gt; Vender (to sell) becomes: vendí: I sold &lt;br&gt; vendiste: you sold &lt;br&gt; vendió: he sold &lt;br&gt; vendimos: we sold &lt;br&gt; vendisteis: you sold &lt;br&gt; vendieron: they sold &lt;br&gt; Conjugate irregular verbs in the preterite. There are 2 types of irregular verbs in the preterite tense. We will start with the more regular of the 2.&lt;br&gt;&lt;br&gt; The endings of these irregulars are all the same, which is why I am referring to them as more regular. The endings are:&lt;br&gt; -e&lt;br&gt; -iste&lt;br&gt; -o&lt;br&gt; -imos&lt;br&gt; -isteis&lt;br&gt; -ieron&lt;br&gt; &lt;br&gt; The stems of these verbs also change as follows:&lt;br&gt; estar --&gt; estuv&lt;br&gt; andar --&gt; anduv&lt;br&gt; tener --&gt; tuv&lt;br&gt; poder --&gt; pud&lt;br&gt; poner --&gt; pus&lt;br&gt; saber --&gt; sup&lt;br&gt; caber --&gt; cup&lt;br&gt; hacer --&gt; hic (exception: 'hizo')&lt;br&gt; venir --&gt; vin&lt;br&gt; querer --&gt; quis&lt;br&gt; decir --&gt; dij&lt;br&gt; traer --&gt; traj&lt;br&gt; and any verb that ends in -ducir --&gt; -duj&lt;br&gt;  Note that for stems ending in j the ending 'ieron' is shortened to 'eron'&lt;br&gt;&lt;br&gt; The more irregular verbs conjugate as follows:&lt;br&gt; "Ir" and "Ser"&lt;br&gt; fui&lt;br&gt; fuiste&lt;br&gt; fue&lt;br&gt; fuimos&lt;br&gt; fuisteis&lt;br&gt; fueron&lt;br&gt;&lt;br&gt; "Dar"&lt;br&gt; di&lt;br&gt; diste&lt;br&gt; dio&lt;br&gt; dimos&lt;br&gt; disteis&lt;br&gt; dieron&lt;br&gt;&lt;br&gt; "Haber"&lt;br&gt; hube&lt;br&gt; hubiste&lt;br&gt; hubo&lt;br&gt; hubimos&lt;br&gt; hubisteis&lt;br&gt; hubieron&lt;br&gt;&lt;br&gt; "Reír"&lt;br&gt; reí&lt;br&gt; reíste&lt;br&gt; rió&lt;br&gt; reímos&lt;br&gt; reísteis&lt;br&gt; rieron&lt;br&gt;&lt;br&gt; "Leer" (similarly Creer)&lt;br&gt; leí&lt;br&gt; leíste&lt;br&gt; leyó&lt;br&gt; leímos&lt;br&gt; leísteis&lt;br&gt; leyeron&lt;br&gt; "Dormir" (similarly Morir)&lt;br&gt; dormí&lt;br&gt; dormiste&lt;br&gt; durmió&lt;br&gt; dormimos&lt;br&gt; dormisteis&lt;br&gt; durmieron&lt;br&gt;&lt;br&gt; "Ver"&lt;br&gt; vi&lt;br&gt; viste&lt;br&gt; vio&lt;br&gt; vimos&lt;br&gt; visteis&lt;br&gt; vieron&lt;br&gt; 

../Verbs/ "He is running. They are dancing. I am liking ... " The above are examples of the present progressive tense, which is used for continuous actions that are taking place right now. This tense works the same way in English and in Spanish. Take the verb "to be" (estar) and conjugate it in the present tense, then add the present participle on the tail end.  Estoy caminando. " I am walking." &lt;br&gt;  Estamos jugando. " We are playing." &lt;br&gt;  Estas fingiendo. " You are faking." &lt;br&gt; 

../Verbs/ Running, jumping, skiing, thinking, talking, ... Corriendo, brincando, esquiando, pensando, hablando ... The present participle is formed by taking off the last two letters of a Spanish verb (-ar, -er, or -ir) and replacing it with -ando (-ar verbs) or -iendo (-er and -ir verbs). Estar &gt; " estando" &lt;br&gt; Comer &gt; " comiendo" &lt;br&gt; Ser &gt; " siendo" &lt;br&gt; Contar &gt; " contando" &lt;br&gt; The slightly irregular ir &gt; "yendo" 

../Verbs/ Broken, thought, gone, ... Roto, pensado, ido, ... The past participle in Spanish is usually the same as the past tense. I thought, I had thought. But sometimes it is the same. I broke, I had broken. In English it is a different conjugation as a rule. Take the -ar, -er, and -ir ending off of the infinitive and replace it with -ado (-ar verbs) or -ido (-er and -ir verbs). Regular conjugation examples. Pensar » "pensado" Comer » "comido" Venir » "venido" Irregular conjugation examples. Romper » "roto" Morir » "muerto" 

Spanish/Verbs I will think, he will shout, you will die, ... Pensaré, gritará, morirás, ... — " or " — Voy a pensar, va a gritar, vas a morir, ... There are two ways to comminicate future events in Spanish. In the first one, add an ending to the unchanged infinitive form of the verb. These same endings are used for all three types of verbs (-AR, -ER and -IR), which makes learning them easier: -é &lt;br&gt; -ás &lt;br&gt; -á &lt;br&gt; -emos &lt;br&gt; -án &lt;br&gt; Pensaré: "I will think" &lt;br&gt; Pensarás: "you will think" &lt;br&gt; Pensará: "it will think" &lt;br&gt; Pensaremos: "we will think" &lt;br&gt; Pensarán: "they will think" &lt;br&gt; Iré: "I will go" &lt;br&gt; Irás: "you will go" &lt;br&gt; Irá: "she will go" &lt;br&gt; Iremos: "we will go" &lt;br&gt; Irán: "y'all will go" &lt;br&gt; There are 12 verbs which change their infinitives before adding the ending, and they can be classified into 3 catigories: First, the "drop 'e's" are "querer," "poder," "caber," "haber," and "saber." Each of these loses the 'e' before the final 'r' when forming the future tense. These can be remembered by the mnemonic "Quick People Can't Have Sushi." Next, the "add 'd's" are "venir," "valer," "salir," "tener," and "poner." For these, the 2nd to last letter of the infinitive is replaced with a 'd'. These can be remembered by the mnemonic "Vroom Vroom, Said The Porsche." Finally, the verbs "decir" and "hacer" change their infinitives to "dir" and "har" respectively before adding the ending. These must simply be memorised. The alternate way to describe the future is to use the present tense of "ir," followed by the infinitive of the action verb. Voy a saber: "I am going to know" &lt;br&gt; Vas a mentir: "You are going to lie" &lt;br&gt; Va a cazar: "He/she is going to hunt" &lt;br&gt; Vamos a contar: "We are going to count" &lt;br&gt; Van a poner: "They are going to put" &lt;br&gt; 

Penetrating oil is useful for removing rust from parts and displacing water. It is more solvent than lubricant and should not be used as a substitute for lubricating oil where lubricating qualities are needed. 

Spanish/Verbs Z. Spanish/Verbs 

The glue used to apply patches to inner tubes. It is a form of contact cement. Unlike most glues, the parts to be joined should be put together after the glue has dried for a few minutes. 

^ Indonesian ^ | « Lesson 8: My Family | Lesson 9: My Home | Lesson 10: At School » 

A pedal can be removed using a pedal wrench or, sometimes, an allen wrench if there is a hexagonal hole in the inside end of the pedal axle. A normal open end spanner/wrench will work on some pedals; an adjustable spanner/wrench spanner will typically not work, because the head is too wide to fit between the pedal and the crank. A good shop-quality pedal wrench will also be longer than these, for the required leverage. This leverage is also the limiting factor for the use of allen wrenches in this application. If a normal spanner is all that you have (be it adjustable or not) then you may find that pedal removal is possible, especially if you are removing platform style pedals—these generally have larger lands for the spanner. Just take care not to scratch or gouge the crank arm (see photo at right). Position an adjustable wrench so that, when turned, the adjustable jaw is on the "inside" of the turn. Usually, you should position the wrench as close to overlapping the crankarm as possible. (So that it is in front of the arm, not extending its line.) Before attempting to remove the right side pedal, make sure the chain is on the largest gear. (That way, if you slip, you'll have a more difficult time impaling yourself on the sharp chainring teeth.) Removing Bicycle Pedals. Note that there are two sizes of pedal axles in common, current-day use: 1/2 inch diameter for the one piece cranks typically found on children's bikes and older American-made bikes, and 9/16 inch diameter for the two and three piece cranks found on most modern adult bikes. Both sizes have SAE threads at 20 tracks per inch. See also. How to Remove and Install Pedals 

=Los deportes (Sports) = 

Thread locking compound is an adhesive designed to prevent parts from vibrating loose. Some brand names are Loc-Tite® and Permatex®. A light duty compound will be removable with simple hand tools. Heavier duty may require special equipment such as a heat gun. There are three standard Loctite grades Loctite was invented in 1953 by Vernon Krieble, PhD. It is described as an anaerobic sealant. In other words, it sets without air. To set it must be in the presence of metal ions and the absence of oxygen. this allows it to set properly inside an assembled fastening. 



oil suitable for use on bicycles can be a basic 30 weight automotive oil or one of the more recent synthetics. Some recent oil products evaporate to a waxy finish on exposed areas, making them less likely to retain grit. 

A standard "pedal wrench" is an open-ended wrench designed to be "thin" enough to fit the narrow wrench flats (gripping surface) typical to pedals. A quality pedal wrench will be "long" enough (and preferably with angled openings) to provide significant leverage, and durable enough to allow repeated application of such force. Pedal wrench flats are "typically 15mm" in size. 9/16" (~14.3mm) is "somewhat common" on older pedals. 17mm and other sizes have been used, but you aren't very likely to encounter them. Pedal wrenches are also available for pedals with "6mm or 8mm internal hex" ("Allen wrench") fittings in the end of the spindle, accessed from the back of the crank arm, through the pedal hole. A "significant portion" of modern pedals provide only this fitting, with no traditional wrench flats. These pedals wrenches differ from standard "Allen wrenches" primarily in that they provide longer and more comfortable handles, which is important to allow application of needed leverage. In the U.S., there are two diameters (sizes) of pedal spindle (axle) threads which you are likely to encounter: 1/2" and 9/16". The first is generally used on bikes with a one piece crank, which will normally be an inexpensive or older bike. The second will usually be found on bicycles with two or three piece cranks, which includes many current multi-speed bicycles and higher quality older bikes. There are other pedal thread sizes, such as 14mm (French), but you are quite unlikely to encounter them in the U.S. WARNING Left-hand pedals are "reverse threaded". (That is, you turn clockwise to "loosen", counterclockwise to "tighten". The Wright brothers introduced this to prevent pedals from unscrewing on their own.) Make very sure it's "not cross-threaded", you are turning the "correct direction", using the "correct pedal", and use "grease"! (If any of the threads are stainless steel or titanium, it would likely be better to use an appropriate "anti-seize" compound instead of grease.) Pedals should be "properly tight", 26 ft-lbs being a generic (not perfect) torque specification. (Rotational force equivalent to 26 pounds on the end of a foot long lever (the wrench), 52 pounds on a 6" lever, etc.) 

Verbs | Verb Tenses | Verbs List El indicativo. These eight simple verb tenses and four compound tenses refer to an objective reality and are categorized by time: present, past, and future. These verb tenses are the simplest and most commonly used. 

../Verbs/ | ../Verb Tenses/ | ../Verbs List/ Subjuntivo. These four verb tenses refer to the speakers' feelings toward the action and are about things that should or could be. These tenses are unfamiliar to English speakers as they are no longer commonly used in English grammar. They are divided into present, past, and future just like the indicative tenses. 

../Verbs/ | ../Verb Tenses/ | ../Verbs List/ El modo potencial. These two verb tenses communicate what would happen or would have happened. 

../Verbs/ | ../Verb Tenses/ | ../Verbs List/ Imperativo. This form gives commands. To form the imperative mood of verbs, one normally takes the present subjunctive of the verb. The major exceptions to this rule are the second person singular and plural forms of positive command. To form a positive second person singular command, use the third person singular indicative. To form a positive second person plural command, take the stem of the verb and add: -ad if the verb is an ar verb. -id if the verb is an ir verb. -ed if the verb is an er verb. Note that when making a negative command, such as "no cantes" (don't sing), the present subjunctive is used. There are many exceptions to these rules. 

../Grammar Index/ 

Vietnamese is a bit different than Romance languages, in that Vietnamese doesn't just use different pronouns for casual or formal situations, but Vietnamese actually uses different pronouns depending on the relation between the speaker and his/her audience. This relation takes gender, age, and status into account. Basically, Vietnamese refer to everyone as a family member. Also, it is common to use a third person personal pronoun in the first person. Singular pronouns. In addition, there are different pronouns for each kind of relative. For a listing of those pronouns, see Family. "(More to come.)" 

= El Departamento de Emergencias (Emergency department) = = Equipo de primeros auxilios (First aid equipment) = 

A degreaser is useful in removing greasy dirt from parts. Traditionally petroleum based products have been used. DO NOT use gasoline. Although it is an effective degreaser, its volatility (i.e. tendency to explode) make it a poor choice for shop work. If you must use a petroleum product, choose naphtha (usually available in hardware stores) or kerosene. A better choice is a modern, biodegradable water based degreaser. Specific blends for bicycle cleaning are available from bike shops and online retailers. You can probably save some money, though, by purchasing degreaser from an auto parts retailer. Whatever you use, after cleaning moving parts with a degreaser, always follow up with lubrication. All residual degreaser should be wiped off the parts before applying lubrication, as it would otherwise dilute the lubricant. 

As with almost all programming languages and Markup Languages there are several different ways of formatting text. HTML is no exception. A popular way to format HTML is to use specific HTML tags such as , and . The World Wide Web consortium (W3C) suggest these tags may not be supported in future and that you use , and to add formatting to text on the HTML page. In practice, these methods are both fine to use, but are limited in the ways that text can be formatted. Thus, many HTML authors are shifting towards using CSS (Cascading Style Sheets) which employs "style sheets" and "inline styles" to add formatting to text. It should be noted that using CSS one can redefine the appearance of text inside all tags of a particular type. See the section on classes below for details. Inline styles. Through CSS (Cascading Style Sheets) you can reuse code using "classes" which contain "styles" (ways of formatting text). You can also put CSS styles inline, as in the following:  This is a paragraph of double-size bold text results in This is a paragraph of double-size bold text Classes. Classes can be defined in two places; in style tags in HTML documents and in external files, the "style sheets" referred to in the name "Cascading Style Sheets". The term "cascading" refers to the mechanism by which the correct style is chosen when multiple style sheets are to be applied to the same document. An example style tag: An example of the contents of a style sheet file: 



About this book An introductory text which should be of particular interest to students of chemical or mechanical engineering. Contents of subject /Introduction/ The basics /Heat Balances/ What goes in must come out /Conduction/ A mechanism of heat transfer /Convection/ A mechanism of heat transfer /Radiation/ A mechanism of heat transfer /Heat Transfer with Phase Change/ Different rules apply when this happens /Heat Exchangers/ Industrial devices for heating, cooling and saving energy /Appendices/ There may be something in here later /Bibliography/ Recommended for more in-depth study 

Welcome. The Voynich manuscript, described as "the world's most mysterious manuscript", is a work which dates to the early 15th century, possibly from northern Italy. It is named after the book dealer /Wilfrid Michael Voynich/, who purchased it in 1912. Some pages are missing, but the current version comprises about 240 pages, most with illustrations. Much of the manuscript resembles herbal manuscripts of the time period, seeming to present illustrations and information about plants and their possible uses for medical purposes. However, most of the plants do not match known species, and the manuscript's script and language remain unknown and unreadable. Possibly some form of encrypted , the Voynich manuscript has been studied by many professional and amateur cryptographers, including American and British codebreakers from both World War I and World War II. As yet, it has defied all decipherment attempts, becoming a "cause célèbre" of historical cryptology. The mystery surrounding it has excited the popular imagination, making the manuscript a subject of both fanciful theories and novels. None of the many speculative solutions proposed over the last hundred years has yet been independently verified. The Voynich manuscript was donated to Yale University's Beinecke Rare Book and Manuscript Library in 1969, where it is catalogued under call number "MS 408" and called a "Cipher Manuscript". Perhaps the appeal of Voynich research is that (a) it is truly cross-disciplinary, and (b) it rewards endeavour and persistence. This wikibook is intended to help you get started on what has already been (for some) a long road of (self-)discovery - the page you may find most useful at first is the Guide To Voynich Jargon. What we know about VMS. Source: What we know about the Voynich manuscript, by Sravana Reddy and Kevin Knight. http://www.isi.edu/natural-language/people/voynich-11.pdf Page-by-page commentary. "See: The Voynich Manuscript/Page-by-page commentary" The Voynich Manuscript has more than 200 pages, divided up into twenty quires (as per the following links). The first 7 quires are standard quires consisting of 4 nested bifolios each. Quire 8 once consisted of 5 such bifolios but only 2 remain. From quire 9 onwards, quires often consist of only one or two multiple folding bifolios. Exceptions are quires 13 (5 standard bifolios) and 20 (7 standard bifolios of which 6 remain). Further details are available at the description of each quire. The web-pages linked onwards from there are intended to summarise the debate relating to individual pages - what could they mean? To what are they similar? What interesting (visual or statistical) properties do they have? "(etc)" Cross-page commentary. This section is designed to contain commentaries on features spanning multiple pages of the VMs. 

This glossary covers only the words used in this textbook. For a more complete listing of Vietnamese words and phrases, please see Wiktionary. 

A multi tool combines several tools into one package.The individual tools may,for example,all fold into the handle. The is an example of a general purpose multi tool. In recent years, a number of manufacturers have produced multi tools specifically intended for bicycle repair. For example: 

Principal Parts. All Latin verbs are identified by four principal parts. By using the four principal parts, one can obtain any and all forms of the verb, including participles, infinitives, gerunds and the like. Examples of principal parts from verbs of each conjugation: For all regular verbs, the principal parts consist of the first person singular present active indicative, the infinitive, the first person singular perfect active indicative, and the supine (or in some texts, the perfect passive participle). Some verbs lack fourth principal parts (e.g., "timeō, timēre, timuī, —"; to be afraid); others, less commonly, lack a third in addition (e.g., "fero, ferre, tuli, latum"; to bring/carry). Others, such as "sum, esse, fuī, futūrus", may use the future active participle ("futūrus") as their fourth principal part; this indicates that the verb cannot be made passive. 

The following people have started lessons as part of this textbook, or have made notable contributions to an existing lesson. They should be able to answer questions you may have about the Vietnamese language or Vietnamese culture: If you feel that you have made notable contributions to this textbook, please feel free to add your name to this list. 

This is a history of the changes that this textbook has gone through: 

Levitov's claim of decipherment has been thoroughly discredited. For a critique by linguistic science, by Jacques B.-M. Guy, see http://www.voynich.net/reeds/gillogly/levitov For a critique by historical analysis, by Dennis J. Stallings, see http://web.archive.org/web/20010410030648/http://www.geocities.com/ctesibos/voynich/levitov2.htm 

Unlike Western languages, Vietnamese doesn't change the ending of the verb (that is, verbs don't conjugate) to express when the statement occurs (the statement's tense). Instead, Vietnamese relies on context to express tense. So, where textbooks for other languages will need to spend chapters and chapters, or pages and pages, to teach you how to change tense, we'll teach you in one page! See what we mean by ? Tense words. One way to do this is to add certain words either before the verb itself, or at the end of the sentence, or both: Present tense. As in most languages, you don't have to do anything special for the present tense – no words expressing tense need to be added. Example: &lt;ruby lang="vi"&gt;&lt;rb&gt;Tôi học&lt;/rb&gt;&lt;br&gt;&lt;rt lang="en-us"&gt;I study&lt;/rt&gt;&lt;rb&gt;&lt;/rb&gt;&lt;/ruby&gt; Present progressive tense. Place "đang" in front of the verb. Example: &lt;ruby lang="vi"&gt;&lt;rb&gt;Tôi&lt;/rb&gt; &lt;rb style="font-weight: bold"&gt;đang&lt;/rb&gt; &lt;rb&gt;học&lt;/rb&gt;&lt;br&gt;&lt;rt lang="en"&gt; I am studying&lt;/rt&gt;&lt;/ruby&gt; Past tense. Either place "đã" in front of the verb, or place "rồi" at the end of the sentence. You can use both, and it often sounds better when you do use both, but it's not necessary, because both words mean "already." Example: &lt;ruby lang="vi"&gt;&lt;rb&gt;Anh ấy&lt;/rb&gt; &lt;rb style="font-weight: bold"&gt;đã&lt;/rb&gt; &lt;rb&gt;học&lt;/rb&gt;&lt;br&gt;&lt;rt lang="en"&gt;He already studied&lt;/rt&gt;&lt;/ruby&gt; can also be written as: &lt;ruby lang="vi"&gt;&lt;rb&gt;Anh ấy học&lt;/rb&gt; &lt;rb style="font-weight: bold"&gt;rồi&lt;/rb&gt;&lt;br&gt;&lt;rt lang="en"&gt;He studied already&lt;/rt&gt;&lt;/ruby&gt; Whereas many other languages have two past tenses – the preterite and imperfect past – Vietnamese has only one. Future tense. Place "sẽ" (will, shall, is about to, plans to) in front of the verb. Example: &lt;ruby lang="vi"&gt;&lt;rb&gt;Anh&lt;/rb&gt; &lt;rb style="font-weight: bold"&gt;sẽ&lt;/rb&gt; &lt;rb&gt;học&lt;/rb&gt;&lt;br&gt;&lt;rt lang="en"&gt;He will study&lt;/rt&gt;&lt;/ruby&gt; Time words. There is another, more natural way of expressing tense. Just use words that express time in the sentence: Present tense. Again, you don't have to do anything special here. But sometimes it's helpful to add hôm nay (today) at the beginning or end of the sentence, just to clear things up. Present progressive tense. Place words like bây giờ (now) or phrases like ngay bây giờ (right now) at the beginning or end of the sentence. Past tense. Place words like hôm qua (yesterday) at the beginning or end of the sentence, or words like mới (recently) right before the verb. You can do both, actually. Future tense. Place words like ngày mai (tomorrow) at the beginning or end of the sentence. Context. Although other methods of expressing tense also rely on context in the sentence, another way to express tense has to do with common sense. For example, in the sentence, . . . "(More to come.)" 

Buoyancy. Buoyancy is the force due to pressure differences on the top and bottom of an object under a fluid (gas or liquid). Net force = buoyant force - force due to gravity on the object Bernoulli's Principle. Fluid flow is a complex phenomenon. An ideal fluid may be described as: As the fluid moves through a pipe of varying cross-section and elevation, the pressure will change along the pipe. The Swiss physicist Daniel Bernoulli (1700-1782) first derived an expression relating the pressure to fluid speed and height. This result is a consequence of conservation of energy and applies to ideal fluids as described above. "Consider an ideal fluid flowing in a pipe of varying cross-section. A fluid in a section of length formula_1 moves to the section of length formula_2 in time formula_3 . The relation given by Bernoulli is: where: In words, the Bernoulli relation may be stated as: "As we move along a streamline the sum of the pressure (formula_4), the kinetic energy per unit volume and the potential energy per unit volume remains a constant." "(To be concluded)" 



Unix and Linux. Overview. Unix is the original Internet Operating System (NFS, TCP/IP, RPC, ...). Linux is a Unix clone written from scratch by Linus Torvalds with assistance from a world wide team of developers across the Net. Properties of both systems are: Linux aims toward POSIX compliance and supports also other standard API's such as BSD and SVR4. Linux has all the features you would expect in a modern Unix, including true multitasking/multiple-users, virtual memory, shared libraries, demand loading, shared copy-on-write executables, proper memory management and TCP/IP networking. More and more companies and organizations are selecting Linux because of its performance, low cost, and royalty free license. Unix and Linux. Unix is the original Internet Operating System written in 1970 developed at A&amp;T Bell Labs. Distributed freely to government and universities. Different entities did support different distributions. Most varieties of Unix carry copyright licenses. Linux was released in 1991 by Linus Torvalds. The model of Linux licence is copy left from GPL. Linux implement most of the System V and BSD Unix commands. Use the man pages to determine the version. Logging In and Out. The logging In session is the first step to access a Linux system. You need a login ID and a password to be able to start using the system. The login ID root is the super user ID that has all privileges, normally assigned to the System Administrator. When a log in session is successful it will set some environment variables and start the shell that has been assigned to the login ID account. Examples of variables set at log in time: $HOME, $SHELL, $PATH, and more. By changing the file /etc/login.defs the administrator can customize some login parameters and variables for all the users. Examples of some of them:  UID_MIN 500  UID_MAX 60000  UMASK 022  ENV_PATH /usr/local/bin:/usr/bin:/bin To logout use the exit or logout command. Password. Nobody can see your password, not even the administrator. To change passwords for users or groups, use passwd.  passwd [options] Common options: As a general rule a password must: User information. For every user account, there is a line defined in the file /etc/passwd. The encrypted password is stored in /etc/shadow. The format of /etc/passwd is:  root:x:0:0:root:/root:/bin/bash  bin:x:1:1:bin:/bin:/bin/bash  daemon:x:2:2:Daemon:/sbin:/bin/bash  yann:x:500:100:Yann Forget:/home/yann:/bin/bash The x field is the old location for the password. The format of /etc/shadow is:  root:IMXweyiV816gg:11996:0:10000::::  bin:!*:8902:0:10000::::  daemon:*:8902:0:10000::::  yann:GoIM8j1S.IuTY:11996:0:99999:7::: The * for the encrypted password means there is no password defined yet. The ! Before the encrypted password means the account is locked. Group information. For every group of users there is a line in file /etc/group. The encrypted password is stored in /etc/gshadow. The format of /etc/group is:  root:x:0:root  bin:x:1:root,bin,daemon  daemon:x:2:  video:x:33:yann The format of /etc/gshadow is:  root:*:root:root  bin:*:root:root  daemon:*:root:root  video:*:root:root,yann Passwords for group are not often implemented since you have to pass to members of a group the password. Misc. Commands. To get information on a logging In session, use id.  id [options] [username] Common options: -g : Print only the group ID. -u : Print only the user ID. Example:  $ id  uid=0(root) gid=0(root) groups=0(root), 1(bin), 14(uucp), 15(shadow),16(dialout) To run a shell with another user ID and group ID, use su.  su [options] [username] Common options: -s : Select another shell. Example:  $ su toto  passwd: Shell. A Shell is a command line interpreter and: There are several popular shells: Most shells let you customize your prompt and environment. For bash the prompt variable is PS1. For bash the file that can be used for the customization of the environment are: Process text streams using filters. Overview. Description: "Candidates should should be able to apply filters to text streams. Tasks include sending text files and output streams through text utility filters to modify the output, and using standard UNIX commands found in the GNU textutils package." Key files terms and utilities include:&lt;br&gt; cat&lt;br&gt; cut&lt;br&gt; expand&lt;br&gt; fmt&lt;br&gt; head&lt;br&gt; join&lt;br&gt; nl&lt;br&gt; od&lt;br&gt; paste&lt;br&gt; pr&lt;br&gt; sed&lt;br&gt; sort&lt;br&gt; split&lt;br&gt; tac&lt;br&gt; tail&lt;br&gt; tr&lt;br&gt; unexpand&lt;br&gt; uniq&lt;br&gt; wc&lt;br&gt; Pattern matching and wildcards. Wildcard is a pattern matching mechanism for filenames. The purpose of wildcard is to increase productivity: Locate files that you don't fully remember. Locate files that have something in common. Work with a group of files rather than individual. The shell interprets these special characters: The characters used for wildcard are: If you use the wildcard characters the shell will try to generate a file from them. Try the following:  echo all files * Special wildcard characters  ? match any one character.  * Any string  [abcfghz] One char set  [a-z] One char in range  [!x-z] Not in set  ~ Home directory  ~user User home directory Examples:  ? One character filenames only  [aA]??? Four characters, starting with a or A.  ~toto pathname of toto home directory  [!0-9]* All string not starting with a number. What about these commands?  ls [a-z][A-Z]??.[uk]  ls big*  ls a???a  ls ??* Shell and wildcards. A shell command line can be a simple command or more complex.  ls -l [fF]*  ls *.c | more  ls -l [a-s]* | mail `users` The first event in the shell is to interpret wildcards. Only the shell interprets unquoted wildcards. Quoting and Comments. Quoting. Do quote to prevent the shell to interpret the special characters and to transform multiple words into one shell word.  echo 'He did it, "Why?"'  echo 'Because "#@&amp;^:-)"'  echo '$VAR='Me  echo "What's happening?"  echo "I don't know but check this $ANSWER"  echo \$VAR=Me  echo What\'s happening\?  echo \\ Comments. You can add comments in a command line or a script. Use the character #. A white space must immediately precede #. Examples:  echo $HOME # Print my Home directory  echo "### PASSED ###" # Only this part is a comment  echo The key h#, not g was pressed. Commands Concatenate files. To concatenate files, use cat.  cat [options] [files...]  tac [options] [files...] The results are displayed to the standard output. Common options:  -s: never more than one single blank line.  -n: number all output line. Examples:  cat file # Display file to the stdout.  cat chapter* # Display all chapters to stdout.  cat -n -s file # Display file with line number with single blank line. To concatenate files in reverse order, use tac. View the beginning and the end of a file. To view only few lines at the beginning or at the end of a file, use head or tail.  head [options] [files...]  tail [options] [files...] The results are displayed to the standard output. Common options:  -n: number of lines to be displayed. (head and tail)  -c: number of bytes to be displayed. (head and tail)  -f: append output. (tail)  -s #: iteration for new data every # sec. (tail) Examples:  head file # Display the first 10 lines of file.  head -n 2 file # Display the first 2 lines of file.  tail -c 10 file # Display the last 10 bytes of file.  tail -f -s 1 /var/log/messages  "Display the last 10 lines of messages, block and check for new data every second." Numbering file lines. To add the line number to a file, use nl.  nl [options] [files...] The results are displayed to the standard output. Common options:  -i #: increment line number by #.  -b: numbering style:  a: number all lines  t: non empty line  n: number no line  -n: numbering format:  rz: right justified  ln: left justified. Examples:  nl file # Add the line number in each line in the file.  nl -b t -n rz file # Add the line number to each non-empty line with zero-completed format. Counting items in a file. To print the number of lines, words and characters of a file, use wc.  wc [options] [files...] The results are displayed to the standard output. Common options:  -c: the byte size.  -m: the number of character.  -w: the number of word.  -l: the number of line.  -L: length of the longest line. Examples:  wc *.[ch] # Display the number of lines, words, and characters for all files .c or .h.  wc -L file # Display the size of the longest line.  wc -w file # Display the number of words. Cutting fields in files. To remove sections from each line of files, use cut.  cut [options] [files...] The results are displayed to the standard output. Common options:  -b #: Extract the byte at position #.  -f #: Extract the field number #. Examples:  cut -b 4 file # Extract and display the 4th byte of each line of file.  cut -b 4,7 file # Extract and display the 4th and 7th byte of each line.  cut -b -2,4-6, 20- file  # Extract character 1 and 2, 4 to 6 and 20 to the end of the line for each line of file.  cut -f 1,3 -d: /etc/passwd # Extract the username and ID of each line in /etc/passwd. The default delimiter is TAB but can be specified with -d. Characters conversion. To translate the stdin to stdout, use tr.  tr [options] SET1 SET2 Common options: -d: delete character in SET1. -s: replace sequence of character in SET1 by one. Examples:  tr ‘a‘ 'A' # Translate lower a with A  tr ‘[A-Z]’ ‘[a-z]’ # Translate uppercase to lowercase  tr -d ‘ ‘ # Delete all spaces from file To convert tabs to spaces, use expand and to convert spaces to tabs, use unexpand.  expand file  unexpand file Lines manipulation. To paste multiple lines of files, use paste.  paste [options] [files...] Common options:  -d #: delimiter: Use # for the delimiter.  -s: serial: paste one file at the time. Examples:  paste f1 f2 # Display line of f1 followed by f2.  paste -d: file1 file2 # Use ':' for the delimiter. To join multiple lines of files, use join.  join file1 file2 To remove duplicated line, use uniq.  uniq [options] [files...] Common options:  -d: only print duplicated lines.  -u: only print unique lines. Examples:  uniq -cd file # Display the number of duplicated line. Splitting files To split big files, use split.  split [options] file Common options:  -l #: split every # lines.  -b #: split file in bytes or b for 512 bytes, k for 1Kbytes, m for 1 Mbytes. Examples:  split -l 25 file # Split file into 25-line files.  split -b 512 file # Split file into 512-byte files.  split -b 2b file # Split file into 2*512-byte files. Formatting for printing. To format a file, use fmt.  fmt [options] [files...] Common options: -w #: maximum line width. Examples:  $ fmt -w 35 file # Display 35-character lines width. To format a file for a printer, use pr.  pr [options] [files...] Common options: -d: double space. Examples:  $ pr -d file # Format file with double-space. Sort lines of text files. To sort the lines of the names files, use sort. sort [options] file The results are displayed to the standard output. Common options:  -r : Reverse  -f : Ignore case  -n : Numeric  -o file: Redirect output to file  -u : No duplicate records  -t; : Use ';' as delimiter, rather than tab or space. Examples:  sort file -r  sort file -ro result Binary file dump. To dump a binary file, use od.  od [options] file The results are displayed to the standard output and start with an offset address in octal. Common options:  -c: each byte as character  -x: 2-byte in hex  -d: 2-byte in decimal  -X: 4-byte in hex.  -D: 4-byte in decimal. Examples:  $ od -cx /bin/ls  0000000 177 E L F 001 001 001 \0 \0 \0 \0 \0 \0 \0 \0 \0  457f 464c 0101 0001 0000 0000 0000 0000  0000020 002 \0 003 \0 001 \0 \0 \0 224 004 \b 4 \0 \0 \0  0002 0003 0001 0000 9420 0804 0034 0000  0000040 ° ² \0 \0 \0 \0 \0 \0 4 \0 \0 006 \0 ( \0  b2b0 0000 0000 0000 0034 0020 0006 0028  0000060 032 \0 031 \0 006 \0 \0 \0 4 \0 \0 \0 4 200 004 \b  001a 0019 0006 0000 0034 0000 8034 0804 Perform basic file management. Overview. Description: "Candidates should be able to use the basic UNIX commands to copy, move, and remove files and directories. Tasks include advanced file management operations such as copying multiple files recursively, removing directories recursively, and moving files that meet a wildcard pattern. This includes using simple and advanced wildcard specifications to refer to files, as well as using find to locate and act on files based on type, size, or time." Key files terms and utilities include:&lt;br&gt; cp&lt;br&gt; find&lt;br&gt; mkdir&lt;br&gt; mv&lt;br&gt; ls&lt;br&gt; rm&lt;br&gt; rmdir&lt;br&gt; touch&lt;br&gt; file globbing&lt;br&gt; cp: Copy files and directories. ls: List directory contents. Create and Remove directories. To create a directory, use mkdir.  mkdir [options] dir Common options:  -m mode: set permission mode. Default use umask.  -p parent: create parent directory as needed. Examples:  mkdir -m 0700 bin  mkdir -p bin/system/x86 To delete an empty directory, use rmdir.  rmdir [options] dir Common options:  -p parent: remove empty subdirectories. Examples:  rmdir tmp  rmdir -p bin/system/x86 Copy files and directories. To copy one file to another, or to a directory, use cp.  cp [options] source target Source and target can be a file or a directory. Common options:  -i interactive: prompt to overwrite  -r recursive: copy the subdirectories and contents. Use -R for special files.  -f force: force the overwriting The default is to silently clobber the target file. Does not alter the source. Examples:  cp *.[a-z] /tmp  cp readme readme.orig  cp ls /bin  cp -ri bin/* /bin Move &amp; Rename files. To rename a file or directory or to move a file or directory to another location, use mv.  mv [options] source target Source and target can be a file or a directory. Common options:  -i interactive: prompt to overwrite  -f force: force the overwriting  -v verbose The default is to silently clobber the target file. Examples:  mv *.[a-z] /tmp  mv readme readme.orig  mv ls /bin  mv -fi bin/* /bin Listing filenames and information. The command to list files in the current directory is ls.  ls [options] [filenames] Common options are:  -l for a long format  -F Append a file type character  -a All files, including hidden  -R Recursive listing of subtree  -d Do not descend into directory The ls is equivalent to the dir command on DOS. Examples of ls output:  $ ls -l /bin/ls  -rwxr-xr-x 1 root root 46784 mar 23 2002 /bin/ls  $ ls -ld /bin  drwxr-xr-x 2 root root 2144 nov 5 11:55 /bin  $ ls -a .  .bash_history .bash_profile .bashrc ...  $ ls -dF /etc .bashrc /bin/ls  .bashrc /bin/ls* /etc/ File types. The long format means:  $ ls -l /etc/hosts #List a long format of the file hosts  -rw-r—r-- 1 root root 677 Jul 5 22:18 /etc/hosts File content and location Linux/Unix does not distinguish file by filename extension, like Windows. To determine the file content use file.  $ file /etc .bashrc /bin/ls /dev/cdrom  /etc: directory  .bashrc: ASCII English text  /bin/ls: ELF 32-bit LSB executable, Intel 80386, version 1 (SYSV), dynamically linked (uses shared libs), stripped  /dev/cdrom: symbolic link to /dev/hdc To determine if a command is a built-in shell command or a program, use type, and use which to find its location.  $ type cp cd ls which type  cp is /bin/cp  cd is a shell builtin  ls is aliased to `ls $LS_OPTIONS'  which is aliased to `type -p'  type is a shell builtin  $ which cut  /usr/bin/cut Creating and using filenames. Filenames can be created with:  cat chapter1 chapter2 &gt; book  vi mynewfile  cp file newfile  netscape  touch memo The valid filename may have:  touch more drink  touch "more drink"  touch memo* Remove files or directories. To remove files or subtree directories, use rm.  rm [options] files Files can be a file or a directory. Common options:  -i interactive: prompt to for each removal  -f force: force the overwriting  -r recursive: remove subtree directories and contents There is no unremove or undelete command. Examples:  rm *.[a-z]  rm readme readme.orig  rm ls /bin  rm -rfi /bin  cd; rm -rf * .* # !! this removes all files of in the home directory of the current user, as well as those in the subdirectories !! Locating files in a subtree directory. To search for a file in a subtree directory, use find.  find [subtrees] [conditions] [actions] The command can take multiple conditions and will search recursively in the subtree. Some possible conditions are:  -name [FNG] # Search fo the FNG name  -type c # Type of file [bcdfl]  -size [+-]# # Has a +- size in blocks (c:bytes,k:kilo)  -user [name] # Own by user  -atime [+-]# # Accessed days ago +n means not been accessed for the last n days -n means been accessed the last ndays.  -mtime [+-]# # Modified days ago  -perm nnn # Has permision flags nnn Some possible actions are:  -print # Print the pathname  -exec cmd {} \; # Execute cmd on the file  -ok cmd {} \; # Same as -exec but ask first Examples  find . -name '*.[ch]' -print  find /var /tmp . -size +20 -print  find ~ -type c -name '*sys*' -print  find / -type f -size +2c -exec rm -i {} \;  find / -atime -3 -print  find ~jo ~toto -user chloe -exec mv {} /tmp \; To locate a binary, source file, or man pages, use whereis.  Whereis [options] Common options:  -b: Search only for binaries.  -m: Search only for manual sections.  -s: Search only for sources. Examples:  $ whereis host  host: /usr/bin/host /etc/host.conf /usr/share/man/man1/host.1.gz  $ whereis -m host  host: /usr/share/man/man1/host.1.gz To locate a file located somewhere defined by the PATH variable, use which.  $ which -a ls  /bin/ls The -a will look for all possible matches in PATH, not just for the first one. Use streams, pipes, and redirects. Overview. Description: "Candidates should be able to redirect streams and connect them in order to efficiently process textual data. Tasks include redirecting standard input, standard output, and standard error, piping the output of one command to the input of another command, using the output of one command as arguments to another command and sending output to both stdout and a file." Key files terms and utilities include:&lt;br&gt; tee&lt;br&gt; xargs&lt;br&gt; «br&gt; «&lt;br&gt; &gt;&lt;br&gt; »&lt;br&gt; Standard input and standard output. For each command executed in a terminal, there are: a standard input value 0. (default keyboard) a standard output value 1.(default terminal) and a standard output for the errors value 2. (default terminal). Each channel can also be identified by an address: &amp;0 for input, &amp;1 for output, And &amp;2 for errors. Each channel [n] can be redirected. [n]&lt; file: Default value of n is 0 and it reads standard input from file. [n]&gt; file: Default value is 1 and it sends standard output to file, overwriting the file if it exists. (clobber) [n]»file: Default value is 1 and it appends standard output to file. «word: Read standard input until word is reached. `command`: Substitute the command name by the result. Examples:  $ pwd &gt; file # out=file in=none error=terminal.  cat chap* &gt;book # out=book in=none error=terminal.  mv /etc/* . 2&gt;error # out=terminal in=none error=error.  echo end of file » book # out=book in=none error=terminal.  set -o noclobber # Shell does not clobber existing files.  ls &gt; list 2&gt;&amp;1 # ls and errors are redirected to list.  ls 2&gt;&amp;1 &gt; list # Errors are redirected to standard output and ls output is redirected to list.  cat `ls /etc/*.conf` &gt; conffile 2»/tmp/errors "Concatenate all the configuration files from /etc dir in conffile and append errors in file /tmp/errors." Redirecting with pipes Pipes are an efficient way to apply multiple commandes concurrently.  command1 | command2 The standard output of command1 will be piped to the standard input of command2. The standard error is not piped. Examples:  ls -l /dev | more  ls -l /etc/*.conf | grep user | grep 500  ls -l /bin | mail `users` To redirect the standard output to a file and to the terminal at the same time, use tee.  ls -l /dev | tee file  ls -l /etc | tee -a file # Append to the file Building arguments The xargs utility constructs an argument list for a command using standard input.  xargs [options] [command] The xargs command creates an argument list for command from standard input. It is typically used with a pipe. Common options: -p: prompt the user before executing each command. Examples:  ls f* | xargs cat # Print to standard output the content of all files starting with f.  find ~ -name 'proj1*' print | xargs cat "Search in the Home directory for files starting with proj1 and send it to the standard input of cat." Use the /dev/null device file to discard output or error messages. Try the following:  grep try /etc/*  grep try /etc/* 2&gt; /dev/null  grep try /etc/* &gt; /dev/null 2&gt; /dev/null Exercises.  ls *.c | xargs rm  ls [aA]* | xargs cat  cat `ls *.v` 2&gt;/dev/null  more `ls *.c` Create, monitor, and kill processes. Overview. Description: "Candidates should be able to manage processes. This includes knowing how to run jobs in the foreground and background, bring a job from the background to the foreground and vice versa, start a process that will run without being connected to a terminal and signal a program to continue running after logout. Tasks also include monitoring active processes, selecting and sorting processes for display, sending signals to processes, killing processes and identifying and killing X applications that did not terminate after the X session closed." Key files terms and utilities include:&lt;br&gt; &amp;&lt;br&gt; bg&lt;br&gt; fg&lt;br&gt; jobs&lt;br&gt; kill&lt;br&gt; nohup&lt;br&gt; ps&lt;br&gt; top Create processes. A running application is a process. Every processes have: A process ID. A parent process ID. A current directory PWD. A file descriptor table. A program which it is executing. Environment variables, inherited from its parent process. Stdin, stdou, and stderr Other Bash is a program that when it is executed becomes a process. Each time you execute a command in a shell a new process is created. Except for the buil-in shell command. They run in the shell context. Use type to check if a command is a built-in shell command. Example type cp ls which type Monitor processes Once the system is up and running from a terminal it is possible to see which processes are running with the ps program. To display a long format of all the processes in the system, do:  ps -Al  F S UID PID PPID C PRI NI ADDR SZ WCHAN TTY TIME CMD  004 S 0 1 0 0 80 0 - 112 do_sel ? 00:00:04 init  004 S 0 381 1 0 80 0 - 332 do_sel ? 00:00:00 dhcpcd  006 S 0 1000 1 0 80 0 - 339 do_sel ? 00:00:00 inetd  044 R 0 1524 1222 0 79 0 - 761 - pts/3 00:00:00 ps The ps program will display all the processes running and their PID numbers and other information. To see a long format of the processes in your login session, do a:  ps -l  F S UID PID PPID C PRI NI ADDR SZ WCHAN TTY TIME CMD  000 S 500 1154 1139 0 80 0 - 724 wait4 pts/1 00:00:00 bash  002 S 500 1285 1283 0 77 0 - 24432 wait_f pts/1 00:00:00 soffice.bin  040 R 500 1442 1435 0 79 0 - 768 - pts/4 00:00:00 ps  F: Process Flags 002: being created, 040: forked but didn't exec, 400: killed by a signal.  S: Process States: R: runnable, S: sleeping, Z: zompbie  UID: User ID, PID: Process ID, PPID: Parent Process ID, C: Scheduler, PRI: priority  NI: Nice value, SZ: size of routine, WCHAN: name of routine Monitor processes. To monitor the processes in real-time, use top. top  9:20am up 2:48, 4 users, load average: 0.15, 0.13, 0.09 78 processes: 75 sleeping, 3 running, 0 zombie, 0 stopped CPU states: 15.3% user, 0.3% system, 0.0% nice, 84.2% idle Mem: 254896K av, 251204K used, 3692K free, 0K shrd, 27384K buff Swap: 514072K av, 0K used, 514072K free 120488K cached  PID USER PRI NI SIZE RSS SHARE STAT %CPU %MEM TIME COMMAND  1517 rarrigon 0 0 40816 39M 17372 R 15.0 16.0 2:59 mozilla-bin  1727 rarrigon 19 0 988 988 768 R 0.3 0.3 0:00 top  1 root 20 0 220 220 188 S 0.0 0.0 0:04 init  2 root 20 0 0 0 0 SW 0.0 0.0 0:00 keventd RSS: The total amount of physical memory used by the task. SHARE: The amount of shared memory used by the task. %CPU: The task's share of the CPU time. %MEM: The task's share of the physical memory. Once top is running it is also possible to execute interactive commands: Type N to sort tasks by pid. Type A to sort tasks by age (newest first). Type P to sort tasks by CPU usage. Type M to sort tasks by memory usage. Type k to kill a process (prompted for pid). Kill processes. The ps program will display all the processes running and their PID numbers. Once the PID is known, it is possible to send signals to the process. SIGSTOP to stop a process. SIGCONT to continue a stopped process. SIGKILL to kill a process. The program to send a signal to a process is called kill.  kill -SIGKILL [pid]  kill -63 [pid]  kill -l By default a process is started in foreground and it is the only one to receive keyboard input. Use CTRL+Z to suspend it. To start a process in background use the &amp;.  bash &amp;  xeyes &amp; Job control In a bash process it is possible to start multiple jobs. The command to manipulate jobs are:  jobs # List all the active jobs  bg %job # Resume job in background  fg %job # Resume job in foreground  kill %job # Kill background job When bash is terminated all processes that have been started from the session will receive the SIGHUP signal. This will by default terminate the process. To prevent the termination of a process, the program can be started with the nohup command.  nohup mydaemon Modify process execution priorities. Overview. Description: "Candidates should should be able to manage process execution priorities. Tasks include running a program with higher or lower priority, determining the priority of a process and changing the priority of a running process." Key files terms and utilities include:&lt;br&gt; nice&lt;br&gt; ps&lt;br&gt; renice&lt;br&gt; top Priorities. To start a command with an adjusted priority, use nice. nice -n +2 [command] nice -n -19 [command] The program nice changes the base time quantum of the scheduler. i.e. it informs the scheduler how important a process is, which is used as a guide to how much CPU time to give it. For example if you wanted to listen to music whilst ripping another CD, you could use nice -n +5 oggenc, you shouldn’t get any “hops” in the music playback, as the scheduler “knows” the oggenc process is less important. The value can go from -19 (highest priority) to +20 (lowest priority). The default value is 0. Only root can set a value below zero. To modify the priority of a running program, use renice. renice +1 -u root # Change the priority for all root processes. renice +2 -p 193 # Change the priority for PID 193 Exercises.  while [ 1 ]  do  echo -n The date is:;  date;  done From the other terminal see that you can stop and continue the print out. 3) Same start as 2) but, make the print out to stop for 3[s] and to continue for 1[s] repeatedly. 4) Make a shell script to renice all process called apache to a 19 value. 5) Do a print from ps formated as: “username”, “command”, “nice value” 6) Kill all the process called “bash” and owned by user polto.. Regular Expression. Overview. Description: "Candidates should be able to manipulate files and text data using regular expressions. This objective includes creating simple regular expressions containing several notational elements. It also includes using regular expression tools to perform searches through a filesystem or file content." Key files terms and utilities include:&lt;br&gt; grep&lt;br&gt; regexp&lt;br&gt; sed Pattern matching. There are two kinds of pattern matching: Wildcard characters are mainly applied when they are used in the current directory or subdirectories. When wildcard characters *, ?, [ - ], ~, and ! are used in regexp they no longer generate filenames. Some of the utilities that use regexp are: Limited regexp search pattern Used by all utilities using regexp. Example:  Ab[0-3]s  ^Ab\^bA  [01]bin$  [^zZ]oro Combinations of limited regexp Combinations used by all utilities using regexp. Examples:  Ab[0-3][a-z]s  ^[01]\^2  [0-9][a-z] \$  [a-zA-Z]*  ^[^c-zC-Z]*  ^[a-zA-Z0-9]$ Modifier patterns Replace strings matched by regexp patterns grep. To find text in a file, use grep.  grep [options] [string] [files] Best to quote the string to prevent interpretation. Common options: -i: Ignore case -l: List filename only if at least one matches -c: Display only count of matched lines -n: Also display line number -v: Must not match. Examples:  grep host /etc/*.conf  grep -l '\&lt;mai' /usr/include/*.h  grep -n toto /etc/group  grep -vc root /etc/passwd  grep '^user' /etc/passwd  grep '[rR].*' /etc/passwd  grep '\&lt;[rR].*' /etc/passwd sed. To apply a command on a stream, use sed.  sed [address1][,address2][!]command[options] [files...] The program sed will apply a command from address1 to address2 in a file. The address1 and address2 format are regular expressions. The sed program is a noninteractive editing tool. Examples:  sed '1,3s/aa/bb/g' file # Replace in file from line 1 to 3 aa by bb.  sed '/here/,$d' file # Delete line from here to the end.  sed '/here/d' file # Delete lines with word here.  sed '1,/xxx/p' file # Print lines from 1 to xxx.  sed '/ll/,/ff/!s/maison/house/g' file # In file replace words maison by house excluding lines from ll to ff. Perform basic file editing using vi. Overview. Description: "Candidates should be able to edit text files using vi. This objective includes vi navigation, basic vi nodes, inserting, editing, deleting, copying, and finding text." Key files terms and utilities include:&lt;br&gt; vi&lt;br&gt; /, ?&lt;br&gt; h,j,k,l&lt;br&gt; G, H, L&lt;br&gt; i, c, d, dd, p, o, a&lt;br&gt; ZZ, :w!, :q!, :e!&lt;br&gt; vi. When using X-Windows, you can use mouse-oriented editors such as xedit. In a cross-development environment, users use their favorite editor. On a non-windowing system, you only need a keyboard editor such as vi. The vi editor on Linux is the same as on any Unix systems. vi has two modes: Transition from one mode to another Command mode to Input mode: i, I, a, A, o, O keys Input mode to Command mode: ESC key The default starting mode is the Command mode The file configuration .exrc can be created in your HOME directory to set up some vi behavior.  set ignorecase # No case-sensitive  set tabs=3 # 3 space for tab character Perform basic file editing using vi Enter Input mode Delete Move cursor 

Create partitions and filesystems. Overview. Description: "Candidates should be able to configure disk partitions and then create filesystems on media such as hard disks. This objective includes using various mkfs commands to set up partitions to various filesystems, including ext2, ext3, reiserfs, vfat, and xfs." Key files terms and utilities include:&lt;br&gt; fdisk&lt;br&gt; mkfs Partitions. Media can be divided into partitions. Partitions are usually created at installation time but can also be created with the fdisk program or other utilities. This will divide the media in partitions and different filesystems can be built and different operating systems can be installed. IDE is recognized as follows: SCSI is recognized as follows: USB and FireWire disks are recognized as SCSI disks. Once it is partitioned it is possible to build a filesystem on every partition. Filesystems. Filesystems exist to allow you to store, retrieve and manipulate data on a media. Filesystems maintain an internal data structure (meta-data) that keeps all you data organized and accessible. The structure of this meta-data gives the characteristic of the filesystem. A filesystem is accessed by a driver through the organized meta-data structure. When Linux boots it reads in /etc/fstab all the filesystems that needs to be mounted and checks if they are in an usable state. When a power failure occurs Linux won't be able to unmout the filesystem properly and some data in the cache won't be synchronized on the media. The meta-data may be corrupted. Once you reboot the system, it will detect this and do a fsck to the full meta-data structure for a consistency check. This can take a very long time. Few minutes to few hours in proportion of the media size. Journaling a filesystem is adding a new data structure called a journal. This journal is on-disk and before each modification to the meta-data is made by the driver it is first written into the journal. Before each meta-data modification the journal maintains a log of the next operation. Now, when a power failure occurs there is only the need to check the journal. A journaling filesystem recovery is very fast. It just needs to go through the logs and fix the latest operation. Journaling filesystem recovery can take only few seconds. On a cluster systems, journaling allows to quickly recover a shared partition of a node that went down. Filesystem tree. A Linux file system has one top directory called root (/) where all sub directories of the entire system is stored. The sub directories can be another partition, a remote directory, or a remote partition accessible through the network with NFS protocol. Create filesystems. To create a file system on a partition, use mkfs.  mkfs [options] -t [fstype] device [blocksize] Common options: The full partition will be erased and organized to the type of filesystem requested. There is no undo command. The fstype possible are: msdos, ext2, ext3, reiserfs,minix,xfs The blocksize allows to customize the block size for your filesystem. Examples:  mkfs -t msdos /dev/fd0  mkfs -t reiserfs /dev/hdd1 4096 Create extended filesystems. To create an extended (ext2, ext3) filesystem on an partition, use mke2fs.  mke2fs [options] device [blocksize] Common options: With mke2fs it is possible to store the super-block the journal information on another device. Examples:  mkefs -b 2048 /dev/fd0 -L floppy  mkfs -V  mke2fs 1.26 (3-Feb-2002) Using EXT2FS Library version 1.263 Monitoring disk usage. To print the disk usage, use du.  du [options] [files...] Common options: Examples:  $ du -ch Documents  112k Documents/Cours/LPI101  4.0k Documents/Cours/LPI102  4.0k Documents/Cours/LPI201  4.0k Documents/Cours/LPI202  124k total  du -sk ~ # Sums up your total disk usage in kilobytes  du -ak ~ | sort -n more # Display every file and its disk space in numerical order. Filesystem disk space. A filesystem is composed of a meta-data structure plus a list of blocks. To print the filesystem disk space usage, use df.  df [options] [files...] Common options: Examples:  $ df -t reiserfs -h  F 1k-blocks Used Available Use% Mounted on  /dev/hda3 28771528 3121536 25649992 11% /  $ df -t ext2 -h  Filesystem Size Used Avail Use% Mounted on  /dev/hda1 15M 3.8M 10M 27% /boot Maintain the integrity of filesystems. Overview. Description: "Candidates should be able to verify the integrity of filesystems, monitor free space and inodes, and repair simple filesystem problems. This objective includes the commands required to maintain a standard filesystem, as well as the extra data associated with a journaling filesystem." Key files terms and utilities include:&lt;br&gt; du&lt;br&gt; df&lt;br&gt; fsck&lt;br&gt; e2fsck&lt;br&gt; mke2fs&lt;br&gt; debugfs&lt;br&gt; dumpe2fs&lt;br&gt; tune2fs Checking filesystems. To check filesystems consistency, use fsck.  fsck [options] -t [fstype] device [fsck-options] Common options: Common fsck-options: Examples:  fsck -t msdos /dev/fd0 -a  fsck -t reiserfs /dev/hda2 -r Checking extended filesystems. To check extended filesystems consistency, use e2fsck.  e2fsck [options] device Common options: Examples:  e2fsck -ay /dev/fd0  e2fsck -f /dev/hda2 Debugging extended filesystems. The debugfs program is an interactive file system debugger. It can be used to examine and change the state of an ext2 file system.  debugfs device Common commands: Example:  stat haut.gif  Inode: 14 Type: regular Mode: 0644 Flags: 0x0 Generation: 67558  User: 0 Group: 0 Size: 3786  File ACL: 0 Directory ACL: 0  Links: 1 Blockcount: 8  Fragment: Address: 0 Number: 0 Size: 0  ctime: 0x3ddf3840 -- Sat Nov 23 09:11:44 2002  atime: 0x3ddf3840 -- Sat Nov 23 09:11:44 2002  mtime: 0x3ddf3840 -- Sat Nov 23 09:11:44 2002  BLOCKS:  (0-3):55-58  TOTAL: 4 Dumping extended filesystems info. To print the super block and blocks group information of an extended filesystem, use dumpe2fs.  dumpe2fs [options] device Common options: Example:  dumpe2fs -h /dev/fd0  dumpe2fs 1.26 (3-Feb-2002)  Filesystem volume name: floppy  Last mounted on: &lt;not available&gt;  Filesystem state: clean  Errors behavior: Continue  Filesystem OS type: Linux  Inode count: 184  Block count: 1440  Reserved block count: 72  Free blocks: 1258  Free inodes: 168  First block: 1  Block size: 1024  First inode: 11  Inode size: 128 Tuning extended filesystems. To tune an extended filesystem, use tune2fs.  tune2fs [options] device Common options: Examples:  tune2fs -L floppy /dev/fd0  tune2fs -l /dev/fd0  (Same output as dumpe2fs -h /dev/fd0)  tune2fs 1.26 (3-Feb-2002)  Filesystem volume name: floppy  Block count: 1440  Reserved block count: 72  Free blocks: 1258  Free inodes: 168  First block: 1  Block size: 1024  First inode: 11  Inode size: 128 Attach and detach filesystems. Overview. Description: "Candidates should be able to configure the mounting of a filesystem. This objective includes the ability to manually mount and unmount filesystems, configure filesystem mounting on bootup, and configure user mountable removable filesystems such as tape drives, floppies, and CDs." Key files terms and utilities include:&lt;br&gt; /etc/fstab&lt;br&gt; mount&lt;br&gt; umount Attach a filesystem. The mount command serves to attach the file system found on some device to the big file tree.  mount [options]  mount [options] [-t vfstype] [-o options] device dir If the device or directory is listed in /etc/fstab you can use the following:  mount [options] [-o options [...]] device | dir Normally only root has the privilege to mount devices unless it is specified in the /etc/fstab file. Examples:  # Print all the mounted filesystems (/etc/mtab).  mount  # Mount devices or dirs listed in /etc/fstab.  mount -a  # Mount /dev/hdc partition in read only mode without updating /etc/mtab.  mount -n -o ro /dev/hdc /mnt  # Allow a user to mount the CDROM if the following line is in /etc/fstab:  # /dev/cdrom /media/cdrom iso9660 ro,user,noauto,unhide  mount /media/cdrom  mount /dev/cdrom  # Sync in realtime  mount -o sync /dev/sdb1 /mnt/usb Detach a filesystem. To detach a filesystem from the file tree, use umount.  umount [options]  umount [options] [-o options [...]] device | dir A busy filesystem cannot be unmounted. Examples:  umount -a # Unmount devices or dirs listed in /etc/fstab.  umount /mnt # Unmount the filesystem attached to /mnt.  umount /media/cdrom # Allow a user to unmount the CDROM if the following line is in /etc/fstab:  /dev/cdrom /media/cdrom iso9660 ro,user,noauto,unhide File system information. The file /etc/fstab contains all the file systems and related information that will be used when doing a mount -a. (Boot time) The file /etc/mtab is maintained by the kernel and keeps track on what is or isn't mounted. The /etc/fstab format is:  #Device Mount point Fs type Options 1 2  /dev/hda3 / reiserfs defaults 1 2  /dev/hda1 /boot ext2 defaults 1 2  /dev/cdrom /media/cdrom auto ro,noauto,user,exec 0 0  usbdevfs /proc/bus/usb usbdevfs noauto 0 0  /dev/hda2 swap swap pri=42 0 0 Common options: Managing disk quota. Overview. Description: "Candidates should be able to manage disk quotas for users. This objective includes setting up a disk quota for a filesystem, editing, checking, and generating user quota reports." Key files terms and utilities include:&lt;br&gt; quota&lt;br&gt; edquota&lt;br&gt; repquota&lt;br&gt; quotaon&lt;br&gt; Quotas. On a system, root can manage the usage of disk space per user and per filsystems. The two limits that can be setup are: The soft limit (soft =) specifies the maximum amount of disk usage a quota user is allowed to have. The hard limit (hard =) specifies the absolute limit on the disk usage a quota user can't go beyond it. There is also the possibility to setup a grace period that will enforce the soft limit only after an amount of time specified. Setting up quotas for users. 1) The keyword usrquota or/and grpquota must be added in file /etc/fstab for the partition interested.  /dev/fd0 /home/yann/mnt auto rw,noauto,user,usrquota 0 0  /dev/hda5 /home ext2 defaults,usrquota,grpquota 1 2 2) Add in each root filesystems the file user.quota or/and group.quota.  touch /mnt/aquota.user  touch /home/aquota.user  touch /home/aquota.group  chmod 600 /mnt/aquota.user  chmod 600 /home/aquota.user  chmod 600 /home/aquota.group Only root can do the quota administration and once the empty files have been created some disk quota can be set such as: Setting up quotas for users. 3) Check the setting  quotacheck -v mnt  quotacheck: Scanning /dev/fd0 [/home/yann/mnt] done  quotacheck: Checked 6 directories and 1 files 4) Enable quota on the disk  quotaon -av  /dev/fd0 [/home/yann/mnt]: user quotas turned on 5) Customize the disk quota limits:  $ edquota -u yann  Disk quotas for user yann (uid 500):  Filesystem locks soft hard inodes soft hard  /dev/fd0 15 0 0 4 0 0  $ edquota -g yann  $ edquota -t  Grace period before enforcing soft limits for users:  Time units may be: days, hours, minutes, or seconds  Filesystem Block grace period Inode grace period  /dev/fd0 7days 7days List quotas. To list quotas for a user or group, use quota.  quota [options] [user|group] Common options: Example:  quota -u yann Display a quota report. To display a quota report, use repquota.  repquota [options] [user|group] Common options: Example:  $ repquota /dev/fd0  *** Report for user quotas on device /dev/fd0  Block grace time: 7days; Inode grace time: 7days  Block limits File limits  User used soft hard grace used soft hard grace  root -- 8 0 0 2 0 0  yann -- 15 0 0 4 0 0 File Permissions and Security. Overview. Unix File Security File and Directory Permissions Default Permissions Changing File Permissions Changing File Owner and Group More Privileges Unix file security The biggest security problem may be you. You own any directory or files you created. You are responsible for accessibility of your files. You decide who can access which files and directories. In your home directory you will be able to grant different levels of permission to yourself, users in your group and the all other users. File and Directory Permissions. The permission of a file or of a directory can be viewed with ls -l. File permissions. Examples of file permissions:  ls -l readme  -rwxrw---- 1 toto users 14 Jul 5 10:00 readme This means read,write, and execution permissions for user toto, read and write permissions for members of group users. No permissions for others. (0760)  ls -l /etc/hosts  -rw-r--r-- 1 root root 14 Jul 5 10:00 /etc/hosts This means read and write permissions for user root, read permissions for members of group root and all others. (0644) Examples of directory permissions:  ls -ld /bin  drwxr-xr-x 2 root root 4096 Jul 5 10:00 /bin This means read,write, and execution permissions for user root, read and execution permissions for members of group root and others. (0755)  ls -l /home/toto  drwxr-xr-x 10 toto users 4096 Jul 5 1:00 /home/toto This means read, write, and execution permissions for user toto, read and execution permissions for members of group users and others. (0755) Default permissions. The default permissions when creating a file are 0666 and when creating a directory are 0777. Most of the systems overwrite this at boot time with the program umask. Generally the mask value is 022. It means the write for group and other will be blocked. To check or change the mask value, do:  umask  umask 066 Examples for file:  default: rw- rw- rw- (0666)  umask: 0 2 2 (0022) Block  result: rw- r-- r-- (0644) Examples for directory:  default: rwx rwx rwx (0777)  umask: 0 2 2 (0022) Block  result: rwx r-x r-x (0755) Changing file permissions. To change permissions on a file or directory, use chmod. To overwrite the existing permissions, do:  chmod 0755 /tmp #rwx for user, rx for group and others To change add or cancel some permissions without overwriting all the existing permissions, do:  chmod u+w readme # Add write permission for user  chmod +r readme # Add read permission for everybody  chmod -r readme # Remove read permission for everybody  chmod u+x,g=r readme # Add execution for user and set read for group  chmod u=rwx,go=rx readme # Set read write and execution for user, read and execution for group and others To change in recursive mode, use the -R option.  chmod -R +x /sbin/* Manage file ownership. Overview. Description: "Candidates should be able to control user and group ownership of files. This objective includes the ability to change the user and group owner of a file as well as the default group owner for new files." Key files terms and utilities include:&lt;br&gt; chmod&lt;br&gt; chown&lt;br&gt; chgrp Changing file owner and group. To change the owner of a file or directory, use chown.  chown yann mon_fichier.txt To change the group of a file or directory, use chgrp.  chgrp dialout caller The programs gpasswd and yast2 allow you to administrate groups. gpasswd [-A user...] [-M user...] group Group administrators can add or delete members of the group  gpasswd -d toto users  gpasswd -a toto users Group administrators can set or remove the password for the group.  gpasswd users  gpasswd -r users More privileges. It is possible to give more privileges to a user when it executes a particular script or program by setting the uid or gid bit of the file. If the bit is set, the process will inherit the permissions of the owner of the file not the permissions of the user. To set the effective uid or gid, use chmod.  chmod 2640 [file] # (2) gid is inheritable for group.  chmod 4640 [file] # (4) uid is inheritable for user. Example of such program is /bin/passwd. The sticky bit can also be set and can make the program text segment resident in RAM. chmod 1640 [file] (1) The file program stays in RAM. Exercises. 1) Write the command line by using letters with chmod to set the following permissions:  rwxrwxr-x :  rwxr--r-- :  r--r----- :  rwxr-xr-x :  rwxr-xr-x :  r-x--x--x :  -w-r----x :  -----xrwx : 2) Write the command line by using octal numbers with chmod to set the following permissions:  rwxrwxrwx :  --x--x--x :  r---w---x :  -w------- :  rw-r----- :  rwx--x--x : 3) With the following umask values what would be the files and directories creation permissions?  umask = 0027  File permissions:  Directory permissions:  umask = 0011  File permissions:  Directory permissions:  umask = 0541  File permissions:  Directory permissions:  umask = 0777  File permissions:  Directory permissions: 4) Create two user accounts Logging in id: tst1, group users, with bash shell, home directory /home/tst1 Logging in id: tst2, group public, with bash shell, home directory /home/tst2 For the two accounts set a password. Logging in as tst1 and copy /bin/ls into tst1 home directory as myls. Change the owner of myls to tst1 and the permissions to 0710. What does this permission value mean? Logging in as tst2 and try to use /home/tst1/myls to list your current directory. Does it work ? Create in /etc/group and /etc/gshadow a new group labo with tst1 and tst2. Change the owner group of myls to labo. Try again from tst2 account to execute /home/tst1/myls to list your current directory. Does it work?. Hard and symbolic links. Overview. Description: "Candidates should be able to create and manage hard and symbolic links to a file. This objective includes the ability to create and identify links, copy files through links, and use linked files to support system administration tasks." Key files terms and utilities include:&lt;br&gt; ln Links. Use link when: You want to create a pathname to a file. Set a shorter or fixed pathname to a file. To link one file to another, use ln:  ln [options] filename linkname  ln [options] filename linkdirectory Common options: The default links are hard links (ln without an option). On Windows they're called shortcuts. A hard link can only be created to an existing file on the same physical device, after creation no visible association can be displayed between a link name and a file name. A symbolic link can be created on a file that doesn’t exist yet, the association between the link name and the file name can be viewed with the ls command. Linking to a file. The symbolic and hard link can be displayed with ls -l. Symbolic link are indicated with an arrow: link_name-&gt;real_filename.  $ ls -l /dev/midi  lrwxrwxrwx 1 root root 6 Jul 4 21:50 /dev/midi -&gt; midi00 Hard links are indicated with the number of links counter (3-1=2 in this case).  $ ls -l readme  -rwxrwxrwx 3 yann users 677 Jul 4 21:50 readme When removing a link name, use rm. Only the link will be removed not the linked file. 

Install &amp; Configure XFree86. Overview. X is a windowing system that provides the basic graphic functions of Linux. It's designed to provide a GUI to any other systems operating across a network, regardless of OS. X operates on a client/server model. X is not part of the operating system. It is an application. The standard Linux X server is Xfree86. The XFree86 log file is located at /var/log/XFree86.0.log. Configuration. type "XFree86 -configure", it will scan your hardware and auto. generate a configuration file matching to your hardware. However, FOR PS/2 MOUSE, I usually need to modify this config file manually from ""Device" "/dev/mouse" to "Device" "/dev/psaux"" Starting and stopping X. To start X you can use:&lt;br&gt; "startx" - terminal command used at level 3;&lt;br&gt; "edit /etc/inittab" to run default at level 5;&lt;br&gt; "xinit" - when there is no .xinitrc file;&lt;br&gt; "init 5" - to change manually runlevel to 5 (and run display manager);&lt;br&gt; "xdm" - (X Display Manager) - graphical login manager, which run automatically at boot process when starting Linux at level 5 (there are also external graphical login managers ex. kdm, gdm).&lt;br&gt; To stop X you can use: &lt;CTRL&gt;+&lt;ALT&gt;+&lt;BACKSPACE&gt;;&lt;br&gt; "init 3" - at lower level than 5 Linux will stop X Window;&lt;br&gt; kill the XFree process. Configuring X To configure X on a system use XF86Setup. The program will generate a configuration file that will be used by the XFree86 server. To tune the screen under X use Xfine2. Under X, the user can configure every conceivable aspect of graphic display. Screen font size, styles Pointer behaviour Screen colors Window manager The tunning can be done on a system-wide or per-user. .xinitrc contains the default window manager and style information to be used by the "startx" command. This file is usually located under /home/username when defined on a per-user basis. .Xdefaults used to setup pointer behaviour, colors, fonts, etc... Setup a display manager. Overview. X needs window managers to manipulate all the graphic applications. Move, size. Open and close. Maximize, minimize, iconize. Title bars The look and feel is implemented in the window manager. Provide virtual desktops. Linux distributions contain many window managers: The desktops bring on top of X Window managers more facilities: Networked X. Overview. X works in a client-server relationship where the application is the client and the server is the application that knows how to provide the services. On a single system, both applications reside on the same system. On a networked system the user can run an X application which is installed on a remote system and do the display on a local system. The advantage of having a X application server is: No local applications installation is needed. No need to have performant local system. No local access to data. Networked X. Exporting a simple application : In order to do so do: startx on the serverhost. xhost + Enable lack of identification on the server host. telnet serverhost from the clienthost and set the DISPLAY variable to the clienthost. Export DISPLAY=clienthost:0.0 Exporting the window manager : In order to do so do: Activate xdmcp with gdmsetup on the server Use Xnest to connect the client on the server in broadcast Exercises. 1)Logging In your system with the failsafe display manager. Execute kde or gdm. Exit kde and logout from failsafe. 2)After logging into your system try to type the following key combinations. CTRL-ALT-F1, CTRL-ALT-F2, CTRL-ALT-F7 What is happening? 3)When you use startx, what is the file where you define the default window manager? 

The following material is directly from the Linux Professional Institute at Objectives 101 reprinted with their permission. Note that their inclusion in this book in no way signifies endorsement on the part of LPI. This is a required exam for LPI certification Level 1. It covers basic system administration skills that are common across all distributions of Linux. Each objective is assigned a weighting value. The weights range roughly from 1 to 10, and indicate the relative importance of each objective. Objectives with higher weights will be covered in the exam with more questions. Topic 101: Hardware &amp; Architecture. 1.101.1 Configure Fundamental BIOS Settings. Description: "Candidates should be able to configure fundamental system hardware by making the correct settings in the system BIOS. This objective includes a proper understanding of BIOS configuration issues such as the use of LBA on IDE hard disks larger than 1024 cylinders, enabling or disabling integrated peripherals, as well as configuring systems with (or without) external peripherals such as keyboards. It also includes the correct setting for IRQ, DMA and I/O addresses for all BIOS administrated ports and settings for error handling." Key files terms and utilities include:&lt;br&gt; /proc/ioports&lt;br&gt; /proc/interrupts&lt;br&gt; /proc/dma&lt;br&gt; /proc/pci&lt;br&gt; 1.101.3 Configure Modem and Sound cards. Description: "Ensure devices meet compatibility requirements (particularly that the modem is NOT a win-modem), verify that both the modem and sound card are using unique and correct IRQ's, I/O, and DMA addresses, if the sound card is PnP install and run sndconfig and isapnp, configure modem for outbound dial-up, configure modem for outbound PPP | SLIP | CSLIP connection, set serial port for 115.2 Kbps." 1.101.4 Setup SCSI Devices. Description: "Candidates should be able to configure SCSI devices using the SCSI BIOS as well as the necessary Linux tools. They also should be able to differentiate between the various types of SCSI. This objective includes manipulating the SCSI BIOS to detect used and available SCSI IDs and setting the correct ID number for different devices especially the boot device. It also includes managing the settings in the computer's BIOS to determine the desired boot sequence if both SCSI and IDE drives are used." Key files terms and utilities include:&lt;br&gt; SCSI ID&lt;br&gt; /proc/scsi/&lt;br&gt; scsi_info&lt;br&gt; 1.101.5 Setup different PC expansion cards. Description: "Candidates should be able to configure various cards for the various expansion slots. They should know the differences between ISA and PCI cards with respect to configuration issues. This objective includes the correct settings of IRQs, DMAs and I/O Ports of the cards, especially to avoid conflicts between devices. It also includes using isapnp if the card is an ISA PnP device." Key files terms and utilities include:&lt;br&gt; /proc/dma&lt;br&gt; /proc/interrupts&lt;br&gt; /proc/ioports&lt;br&gt; /proc/pci&lt;br&gt; pnpdump(8)&lt;br&gt; isapnp(8)&lt;br&gt; lspci(8)&lt;br&gt; 1.101.6 Configure Communication Devices. Description: "Candidates should be able to install and configure different internal and external communication devices like modems, ISDN adapters, and DSL switches. This objective includes verification of compatibility requirements (especially important if that modem is a winmodem), necessary hardware settings for internal devices (IRQs, DMAs, I/O ports), and loading and configuring suitable device drivers. It also includes communication device and interface configuration requirements, such as the right serial port for 115.2 Kbps, and the correct modem settings for outbound PPP connection(s)." Key files terms and utilities include:&lt;br&gt; /proc/dma&lt;br&gt; /proc/interrupts&lt;br&gt; /proc/ioports&lt;br&gt; setserial(8)&lt;br&gt; 1.101.7 Configure USB devices. Description: "Candidates should be able to activate USB support, use and configure different USB devices. This objective includes the correct selection of the USB chipset and the corresponding module. It also includes the knowledge of the basic architecture of the layer model of USB as well as the different modules used in the different layers.." Key files terms and utilities include:&lt;br&gt; lspci(8)&lt;br&gt; usb-uhci.o&lt;br&gt; usb-ohci.o&lt;br&gt; /etc/usbmgr/&lt;br&gt; usbmodules&lt;br&gt; /etc/hotplug&lt;br&gt; Topic 102: Linux Installation &amp; Package Management. 1.102.1 Design hard disk layout. Description: "Candidates should be able to design a disk partitioning scheme for a Linux system. This objective includes allocating filesystems or swap space to separate partitions or disks, and tailoring the design to the intended use of the system. It also includes placing /boot on a partition that conforms with the BIOS' requirements for booting." Key files terms and utilities include:&lt;br&gt; / (root) filesystem&lt;br&gt; /var filesystem&lt;br&gt; /home filesystem&lt;br&gt; swap space&lt;br&gt; mount points&lt;br&gt; partitions&lt;br&gt; cylinder 1024&lt;br&gt; 1.102.2 Install a boot manager. Description: "Candidate should be able to select, install, and configure a boot manager. This objective includes providing alternative boot locations and backup boot options (for example, using a boot floppy)." Key files terms and utilities include:&lt;br&gt; /etc/lilo.conf&lt;br&gt; /boot/grub/grub.conf&lt;br&gt; lilo&lt;br&gt; grub-install&lt;br&gt; MBR&lt;br&gt; superblock&lt;br&gt; first stage boot loader&lt;br&gt; 1.102.3 Make and install programs from source. Description: "Candidates should be able to build and install an executable program from source. This objective includes being able to unpack a file of sources. Candidates should be able to make simple customizations to the Makefile, for example changing paths or adding extra include directories." Key files terms and utilities include:&lt;br&gt; gunzip&lt;br&gt; gzip&lt;br&gt; bzip2&lt;br&gt; tar&lt;br&gt; configure&lt;br&gt; make&lt;br&gt; 1.102.4 Manage shared libraries. Description: "Candidates should be able to determine the shared libraries that executable programs depend on and install them when necessary. Candidates should be able to state where system libraries are kept." Key files terms and utilities include:&lt;br&gt; ldd&lt;br&gt; ldconfig&lt;br&gt; /etc/ld.so.conf&lt;br&gt; LD_LIBRARY_PATH&lt;br&gt; 1.102.5 Use Debian package management. Description: "Candidates should be able to perform package management skills using the Debian package manager. This objective includes being able to use command-line and interactive tools to install, upgrade, or uninstall packages, as well as find packages containing specific files or software (such packages might or might not be installed). This objective also includes being able to obtain package information like version, content, dependencies, package integrity and installation status (whether or not the package is installed)." Key files terms and utilities include:&lt;br&gt; unpack&lt;br&gt; configure&lt;br&gt; /etc/dpkg/dpkg.cfg&lt;br&gt; /var/lib/dpkg/*&lt;br&gt; /etc/apt/apt.conf&lt;br&gt; /etc/apt/sources.list&lt;br&gt; dpkg&lt;br&gt; dselect&lt;br&gt; dpkg-reconfigure&lt;br&gt; apt-get&lt;br&gt; alien&lt;br&gt; 1.102.6 Use Red Hat Package Manager (RPM). Description: "Candidates should be able to perform package management under Linux distributions that use RPMs for package distribution. This objective includes being able to install, re-install, upgrade, and remove packages, as well as obtain status and version information on packages. This objective also includes obtaining package information such as version, status, dependencies, integrity, and signatures. Candidates should be able to determine what files a package provides, as well as find which package a specific file comes from." Key files terms and utilities include:&lt;br&gt; /etc/rpmrc&lt;br&gt; /usr/lib/rpm/*&lt;br&gt; rpm&lt;br&gt; grep&lt;br&gt; Topic: 103 GNU &amp; Unix Commands. 1.103.1 Work on the command line. Description: "Candidates should be able to Interact with shells and commands using the command line. This includes typing valid commands and command sequences, defining, referencing and exporting environment variables, using command history and editing facilities, invoking commands in the path and outside the path, using command substitution, applying commands recursively through a directory tree and using man to find out about commands." Key files terms and utilities include:&lt;br&gt; .&lt;br&gt; bash&lt;br&gt; echo&lt;br&gt; env&lt;br&gt; exec&lt;br&gt; export&lt;br&gt; man&lt;br&gt; pwd&lt;br&gt; set&lt;br&gt; unset&lt;br&gt; ~/.bash_history&lt;br&gt; ~/.profile&lt;br&gt; 1.103.2 Process text streams using filters. Description: "Candidates should should be able to apply filters to text streams. Tasks include sending text files and output streams through text utility filters to modify the output, and using standard UNIX commands found in the GNU textutils package." Key files terms and utilities include:&lt;br&gt; cat&lt;br&gt; cut&lt;br&gt; expand&lt;br&gt; fmt&lt;br&gt; head&lt;br&gt; join&lt;br&gt; nl&lt;br&gt; od&lt;br&gt; paste&lt;br&gt; pr&lt;br&gt; sed&lt;br&gt; sort&lt;br&gt; split&lt;br&gt; tac&lt;br&gt; tail&lt;br&gt; tr&lt;br&gt; unexpand&lt;br&gt; uniq&lt;br&gt; wc&lt;br&gt; 1.103.3 Perform basic file management. Description: "Candidates should be able to use the basic UNIX commands to copy, move, and remove files and directories. Tasks include advanced file management operations such as copying multiple files recursively, removing directories recursively, and moving files that meet a wildcard pattern. This includes using simple and advanced wildcard specifications to refer to files, as well as using find to locate and act on files based on type, size, or time." Key files terms and utilities include:&lt;br&gt; cp&lt;br&gt; find&lt;br&gt; mkdir&lt;br&gt; mv&lt;br&gt; ls&lt;br&gt; rm&lt;br&gt; rmdir&lt;br&gt; touch&lt;br&gt; file globbing&lt;br&gt; 1.103.4 Use streams, pipes, and redirects. Description: "Candidates should be able to redirect streams and connect them in order to efficiently process textual data. Tasks include redirecting standard input, standard output, and standard error, piping the output of one command to the input of another command, using the output of one command as arguments to another command and sending output to both stdout and a file." Key files terms and utilities include:&lt;br&gt; tee&lt;br&gt; xargs&lt;br&gt; «br&gt; «&lt;br&gt; &gt;&lt;br&gt; »&lt;br&gt; ` `&lt;br&gt; 1.103.5 Create, monitor, and kill processes. Description: "Candidates should be able to manage processes. This includes knowing how to run jobs in the foreground and background, bring a job from the background to the foreground and vice versa, start a process that will run without being connected to a terminal and signal a program to continue running after logout. Tasks also include monitoring active processes, selecting and sorting processes for display, sending signals to processes, killing processes and identifying and killing X applications that did not terminate after the X session closed." Key files terms and utilities include:&lt;br&gt; &amp;&lt;br&gt; bg&lt;br&gt; fg&lt;br&gt; jobs&lt;br&gt; kill&lt;br&gt; nohup&lt;br&gt; ps&lt;br&gt; top&lt;br&gt; 1.103.6 Modify process execution priorities. Description: "Candidates should should be able to manage process execution priorities. Tasks include running a program with higher or lower priority, determining the priority of a process and changing the priority of a running process." Key files terms and utilities include:&lt;br&gt; nice&lt;br&gt; ps&lt;br&gt; renice&lt;br&gt; top&lt;br&gt; 1.103.7 Search text files using regular expressions. Description: "Candidates should be able to manipulate files and text data using regular expressions. This objective includes creating simple regular expressions containing several notational elements. It also includes using regular expression tools to perform searches through a filesystem or file content." Key files terms and utilities include:&lt;br&gt; grep&lt;br&gt; regexp&lt;br&gt; sed&lt;br&gt; 1.103.8 Perform basic file editing operations using vi. Description: "Candidates should be able to edit text files using vi. This objective includes vi navigation, basic vi nodes, inserting, editing, deleting, copying, and finding text." Key files terms and utilities include:&lt;br&gt; vi&lt;br&gt; /, ?&lt;br&gt; h,j,k,l&lt;br&gt; G, H, L&lt;br&gt; i, c, d, dd, p, o, a&lt;br&gt; ZZ, :w!, :q!, :e!&lt;br&gt; Topic 104: Devices, Linux Filesystems, Filesystem Hierarchy Standard. 1.104.1 Create partitions and filesystems. Description: "Candidates should be able to configure disk partitions and then create filesystems on media such as hard disks. This objective includes using various mkfs commands to set up partitions to various filesystems, including ext2, ext3, reiserfs, vfat, and xfs." Key files terms and utilities include:&lt;br&gt; fdisk&lt;br&gt; mkfs&lt;br&gt; 1.104.2 Maintain the integrity of filesystems. Description: "Candidates should be able to verify the integrity of filesystems, monitor free space and inodes, and repair simple filesystem problems. This objective includes the commands required to maintain a standard filesystem, as well as the extra data associated with a journaling filesystem." Key files terms and utilities include:&lt;br&gt; du&lt;br&gt; df&lt;br&gt; fsck&lt;br&gt; e2fsck&lt;br&gt; mke2fs&lt;br&gt; debugfs&lt;br&gt; dumpe2fs&lt;br&gt; tune2fs&lt;br&gt; 1.104.3 Control mounting and unmounting filesystems. Description: "Candidates should be able to configure the mounting of a filesystem. This objective includes the ability to manually mount and unmount filesystems, configure filesystem mounting on bootup, and configure user mountable removeable filesystems such as tape drives, floppies, and CDs." Key files terms and utilities include:&lt;br&gt; /etc/fstab&lt;br&gt; mount&lt;br&gt; umount&lt;br&gt; 1.104.4 Managing disk quota. Description: "Candidates should be able to manage disk quotas for users. This objective includes setting up a disk quota for a filesystem, editing, checking, and generating user quota reports." Key files terms and utilities include:&lt;br&gt; quota&lt;br&gt; edquota&lt;br&gt; repquota&lt;br&gt; quotaon&lt;br&gt; 1.104.5 Use file permissions to control access to files. Description: "Candidates should be able to control file access through permissions. This objective includes access permissions on regular and special files as well as directories. Also included are access modes such as suid, sgid, and the sticky bit, the use of the group field to grant file access to workgroups, the immutable flag, and the default file creation mode." Key files terms and utilities include:&lt;br&gt; chmod&lt;br&gt; umask&lt;br&gt; chattr&lt;br&gt; 1.104.6 Manage file ownership. Description: "Candidates should be able to control user and group ownership of files. This objective includes the ability to change the user and group owner of a file as well as the default group owner for new files." Key files terms and utilities include:&lt;br&gt; chmod&lt;br&gt; chown&lt;br&gt; chgrp&lt;br&gt; 1.104.7 Create and change hard and symbolic links. Description: "Candidates should be able to create and manage hard and symbolic links to a file. This objective includes the ability to create and identify links, copy files through links, and use linked files to support system administration tasks." Key files terms and utilities include:&lt;br&gt; ln&lt;br&gt; 1.104.8 Find system files and place files in the correct location. Description: "Candidates should be thoroughly familiar with the Filesystem Hierarchy Standard, including typical file locations and directory classifications. This objective includes the ability to find files and commands on a Linux system." Key files terms and utilities include:&lt;br&gt; find&lt;br&gt; locate&lt;br&gt; slocate&lt;br&gt; updatedb&lt;br&gt; whereis&lt;br&gt; which&lt;br&gt; /etc/updatedb.conf&lt;br&gt; Topic 110: The X Window System. 1.110.1 Install &amp; Configure XFree86. Description: "Candidate should be able to configure and install X and an X font server. This objective includes verifying that the video card and monitor are supported by an X server, as well as customizing and tuning X for the videocard and monitor. It also includes installing an X font server, installing fonts, and configuring X to use the font server (may require a manual edit of /etc/X11/XF86Config in the "Files" section)." Key files terms and utilities include:&lt;br&gt; XF86Setup&lt;br&gt; xf86config&lt;br&gt; xvidtune&lt;br&gt; /etc/X11/XF86Config&lt;br&gt; .Xresources&lt;br&gt; 1.110.2 Setup a display manager. Description: "Candidate should be able setup and customize a Display manager. This objective includes turning the display manager on or off and changing the display manager greeting. This objective includes changing default bitplanes for the display manager. It also includes configuring display managers for use by X-stations. This objective covers the display managers XDM (X Display Manger), GDM (Gnome Display Manager) and KDM (KDE Display Manager)." Key files terms and utilities include:&lt;br&gt; /etc/inittab&lt;br&gt; /etc/X11/xdm/*&lt;br&gt; /etc/X11/kdm/*&lt;br&gt; /etc/X11/gdm/*&lt;br&gt; 1.110.4 Install &amp; Customize a Window Manager Environment. Description: "Candidate should be able to customize a system-wide desktop environment and/or window manager, to demonstrate an understanding of customization procedures for window manager menus and/or desktop panel menus. This objective includes selecting and configuring the desired x-terminal (xterm, rxvt, aterm etc.), verifying and resolving library dependency issues for X applications, exporting X-display to a client workstation." Key files terms and utilities include:&lt;br&gt; .xinitrc&lt;br&gt; .Xdefaults&lt;br&gt; xhost&lt;br&gt; DISPLAY environment variable&lt;br&gt; 

Chinese, like all languages, has its own set of unique greetings which may be seemingly strange to learners of the language (this is particularly true if the two cultures are vastly different). Below, you will find commonly-used Mandarin greetings and farewells, along with corresponding pinyin pronunciations. 

= El amor = 

=Lesson 2: 今天你忙不忙？= Lesson 2 contains a dialogue of two students discussing their classes for the day. Dialogues. Dialogue 1 Dialogue 2 Vocabulary. Note: Visit this lesson's Stroke Order subpage to see images and animations detailing how to write the following characters. Audio files of the words are linked from the pīnyīn when available. Problems listening? See . Grammar. The adverb Hěn [很]. &lt;br&gt; 1. 我很忙。 Le [了] as emphasizer. &lt;br&gt; 1. 太多了。 2. 太少了。 Affirmative-negative questions. &lt;br&gt; &lt;br&gt;&lt;br&gt; Example:&lt;br&gt; Q: 他是不是东尼？ A: 是的。（是，他是/嗯，他是。）or 不是。 （不，他不是。） &lt;br&gt; &lt;br&gt;&lt;br&gt; Example: &lt;br&gt; Q:艾美今天忙不忙？/艾美今天忙不？ A: 她很忙。or 她不忙。 Sentences using yǒu [有]. &lt;br&gt; &lt;br&gt;&lt;br&gt; Example: &lt;br&gt; 我有三门课。 &lt;br&gt; &lt;br&gt;&lt;br&gt; Example: &lt;br&gt; 今天，他们没有课。 &lt;br&gt; &lt;br&gt;&lt;br&gt; Example: &lt;br&gt; 今天，他们一门课都没有。 

= Lesson 3: 助詞 = The Chinese language employs heavy usage of particles to modify the meaning of characters and sentences. Since Chinese has neither inflections nor tense, the mastery of particles is an absolute must if one is to fully comprehend both written and spoken Chinese. Below, you will find some of the most common particles in everyday Chinese. The De {的} particle. Example:  她 的 名字 是 金妮。  Tā de míngzi shì Jīnní.  Her name is Ginny. Example 她是一个美丽的姑娘  Tā shì yīge měilì de gū’niang.  She is a beautiful girl. where "美丽" "beautiful" is an adjective, and Example 研究是一个科学的过程  Yánjīu shì yígè kēxué dè guòchéng  Researching is a scientific process. and where "科学" is a noun in Chinese and is turned into adjective using "的". The Le {了} / Liăo {了} particle. Perfective Aspect Particle Example:  他 走 了。  Tā zŏu le.  He has gone. ※The "le" here is used to modify 走 (zŏu, "to go") into an action which has already been completed. Example:  別 再 打扰 我 了!  別 再 打擾 我 了!  Bié zài dărăo wŏ le!  Do not bother me again! ※In this instance, le is used in conjunction with bié ("do not") to form an imperative. "Note": most imperatives are not formed using this construction. Example:  我 实在 吃 不 了 了。  我 實在 吃 不 了 了。  Wŏ shízài chī bù liăo le.  I cannot possibly eat any more. At first glance, this sentence may seem a bit daunting as it includes two instances of the le particle, paired side-by-side. However, the first le is understood to be liăo given its placement (bù + le is a nonsensical pairing). Therefore, liăo serves to indicate the capability of eating any further and le "emphasizes" this assertion. The Zhe [着] particle showing continuation. 1. 他睡着觉时有人敲门。 The Zháo [着] particle indicating accomplishment. 1. 我终于把东西买着了!  (我終於把東西買著了!) And another word, dào [到], can be seen as a substitution for 着, in most cases they are interchangeable. 2. 他在行窃时被当场抓到。 The 把 + N + V + 着(到)了 construction is particularly useful and should be studied. The De [得] particle indicating degree. 1. 我说得很好. This construct often requires a context to gain its full meaning. If you wish to speak more specifically about an action, the two constructs below demonstrate the use of 得 with a direct object. 2. 我说中文说得很好. Note the dual-use of the verb. 3. 中文我说得很好. This construct emphasizes the object (here being "Chinese"). 4. 我中文说得很好. This expression is the simplification of the 2nd expression by eliminating the 1st verb. This form is even more frequently used than the 2nd expression above. Vocabulary. Note: Visit this lesson's Stroke Order subpage to see images and animations detailing how to write the following characters. Audio files of the words are linked from the pīnyīn when available. Problems listening? See . 

Synopsis. This book attempts to teach the skills that can help you to have lucid dreams — dreams in which you know that you are dreaming. For the skeptical, lucid dreams have been scientifically demonstrated to exist. The ability to lucid dream will open your mind to a world of infinite possibilities as you become adept at taking control of your dreams. We will start by explaining how lucid dreaming works biologically. Next, the book will prepare you for lucid dreaming by helping you to remember more of your dreams (dream recall). You will then learn a variety of ways of becoming, and staying, lucid. Finally, you will find suggested activities to try while in the dream world. Contents. Before each target there is an image with a subjective indication of how complete that target is: 

&lt;br&gt;&lt;br&gt; =Lesson 4: Word order and Verbs= Basic Word Order. Subject-Verb-Object. The order of most Chinese sentences, like in English, is S-V-O, that is Subject-Verb-Object. Word order in Chinese is more rigid than in English. However, sometimes you may find sentences that seem to defy normal word order. For example, 我住在中国。wǒ zhù zài zhōngguó. The English translation does this too: I live in China. The reason for this is that "in China" is a preposition (prepositions indicate place or time) that is tacked on to the main sentence—"I live." More examples: As in English, a preposition can also appear after a subject. When using both a preposition for time and a preposition for place, put the preposition for time first. Note the variation in word order. You can also place a preposition for place, but not for time, at the end of a sentence. Topic-Comment. Another structure for Chinese sentences is topic-comment. That is, the first thing mentioned is the topic of discussion and then the speaker will add a comment following that. It is used to emphasize a certain part of the sentence. In the following example, the speaker wants to emphasize that he is going to read the particular book being discussed. Zhè běn shū, wǒ míngtiān zài wǒ jiā kàn. 

Content needed. AP (Advanced Placement) courses allow high school students to take and get credit for college-level classes. To get credit for a course, a student must pass the AP test. Unfortunately, not all colleges and universities in the US accept AP tests, and some may have their own independent exams. The AP is prepared by Collegeboard and their buddies at ETS (I think), also the writers of such tests as the SAT and PSAT. Courses cover a wide range of subjects, including Spanish, Calculus, English, Chemistry, and History. 1. Agree or Disagree&lt;br&gt; 2. Author's attitude, tone, or point of view&lt;br&gt; 3. Analyze effect&lt;br&gt; 4. Author's purpose&lt;br&gt; 5. General statement about society or human nature&lt;br&gt; 6. Comparison/Contrast&lt;br&gt; 

Chapter 10 Chapter 10. Mycology ~ The Fungi Introduction. The Fungi (singular is "fungus") are a large group of organisms treated within the science of Botany, but not really "plants" in the usual sense of the term. The Fungi are ranked as a kingdom within the Domain . That is, they have eukaryotic cells with distinct nuclei, although in some species divisions between nucleated "cells" are sparse. However, they all lack chlorophyll, and the species are saprophytic, parasitic or mycorrhizal. Included within the Fungi are the well known mushrooms, but the group also includes many microscopic forms, and fungi inhabit every environment on earth, perhaps second only to the bacteria (Chapter 7) in distribution. 

This page describes a number of lucid dream induction techniques. It is recommended that you be able to recall "at least" one dream per night in order to maximize the effectiveness of these methods. Preliminary Knowledge. Certain elements are common to many of the lucidity-inducing techniques discussed later in this chapter. To better understand these techniques, these common components will be discussed first. Sleep Interruption. An element shared by many of the techniques is sleep interruption. Sleep interruption is the process of purposefully awakening during your normal sleep period and falling asleep a short time later (10–60 minutes). This can be easily done by using a relatively quiet alarm clock to bring you to consciousness without fully waking you. If you find yourself resetting the clock in your sleep, it can be placed on the other side of the room, forcing you to get out of bed to turn it off. Other biorhythm-based options involve drinking lots of fluid (particularly water or tea, a known diuretic) prior to sleep, forcing one to get up to urinate. Sleep interruption is a natural part of the MILD technique (described below) which trains you to arise immediately after your dreams end. Sleep Continuity. If you have trouble initially falling asleep, avoid drinking water for about an hour before going to bed. Otherwise, you may find yourself running to the bathroom, disrupting any attempts at lucidity. Also, try to avoid caffeine and sugar before bed. However, depending on your sensitivity, caffeine may only stimulate your mind as opposed to your body. This extra grip on consciousness could be helpful in inducing lucid dreams. Exercising during the day is an excellent way of preparing your body for sleep. However, be sure to not exercise inside the three hours before bedtime, as your body will be stimulated for a short time afterwards. The morning or afternoon is the best time for this. If you still have difficulty getting to sleep, try reading about lucid dreaming just before going to sleep. Your subconscious will likely absorb this information, increasing your chances of experiencing a lucid dream. If you "do" decide to read before going to sleep, keep a lamp next to your bed as physically getting up to turn off the lights may reawaken your body. Reality checks. A reality check is a test you can perform to see if you're dreaming or awake. It might seem odd to test reality when you are sure that you're awake, but making a habit out of one or more of these reality checks will hugely increase your chances of having a lucid dream. If, say, you hold your nose and try to breathe in through it several times throughout the day then you're very likely to dream about doing it. And when you dream about performing a reality check, then of course the results should come out differently, in this case you'll find that you are somehow breathing in through your closed nostrils. You'll know that you're dreaming, and be able to take lucid control! So here are some reality checks. You should be familiar with the entire list even if you only use a few. Choose a few reality checks which you will do regularly. Take them seriously, do not assume you are awake (even when you know you are). If you practice performing these checks very thoroughly while awake, then you're more likely to perform them thoroughly while dreaming. You should always carry out more than one reality check. If you find that it is not a dream, look around you and think of what would be different if it were a dream. If you do this it will make it more likely that you will do a reality check in a dream. Apart from doing reality checks throughout the day, you also need to do a reality check immediately after you wake up. This helps you become lucid in false awakenings, when you dream that you have woken up but in fact are still dozing. Using a digital alarm clock or mobile phone display to do a Reading check, every single time you wake up, is a quick and reliable way to catch false awakenings. If you have trouble bringing reality checks into your dreams then before going to bed imagine yourself in a dream, noticing odd details and doing a reality check. Then do a reality check in real life. If you do this a few times before bed you will find that you will do it more often in dreams. If you are in a situation where you cannot do a reality check, such as at a public speaking, try to do one as soon as possible. You can do some reality checks very inconspicuously, such as looking at some text on a sign. If you start to say “well, I can't do a reality check now” you should not be surprised when you make this mistake in a dream! Which reality checks are best? When selecting reality checks, the most important properties of a reality check are "reliability", "speed", and "inconspicuousness". On the table above, these are scored out of 5. I have trouble remembering to do reality checks throughout the day. What reminders can I use? You are lucky to have an interesting day and forget about lucid dreaming! You can tag your mind to remember dreaming when you think of certain things, like your friend or your homework. It isn't advisable to explicitly write “reality check” or “lucid” on your hand, as this could create an over dependence on this reminder, which may not exist in a dream. However, you might want to just draw a dot or small circle on your hand. This should be enough to remind you to do a reality check. Try putting a little label on your clock, mobile phone, or watch, reminding yourself to do a reality check. (Some weird colors will make it more noticeable and it will take longer for you to get used to it and ignore it.) If you check these regularly during the course of your waking day, you will be doing lots of reality checks. A simple coffee mug with a reminder such as "Are you dreaming?" printed on it or random alarms can also serve well, but try not to become too dependent on them. You can find examples of these at LD4All. Another technique is to write down three things you do regularly in a day. Examples include hearing your name, going through a doorway, turning on a TV, beginning to read a book, or seeing a stranger. In the morning, choose three such events and intend to do a reality check whenever they happen in the following day. I did a reality check in a dream but it said that I was not dreaming. What went wrong? Some reality checks work perfectly for some people and awfully for others. These are mostly the light switches one and the hands one. If you find that the light switch works or that your hands are perfectly normal, you need to change to a different technique. I did a reality check in a dream but I did not quite realize I was dreaming. What went wrong? An example of this is looking into a mirror and seeing some huge boils or a gray mist on your reflection and not realizing that you are dreaming. This is rare if you actually intended to look into the mirror as a reality check. You need to be more careful when doing your reality checks in real life or pick more reliable reality checks which show more obviously that you are dreaming. Also try to pick reality checks that are easy to do. For example, don't rely the Time Reality Check if you never wear a watch, and don't pick the Mirror Reality Check if you rarely look in the mirror during the day or you know that you won't find a mirror in your dream. Another good remedy for this problem (and a good practice in general) is to perform two or three reality checks at a time. The Time Reality Check, for example, can be easily combined with attempting to push one hand through another. Or, for those with glasses, testing your ability to read text fits naturally with checking for "perfect eyesight w/o glasses". Threads about reality checks on ld4all.com: The BIG Reality Check topic I II | RCs that prove you are awake while dreaming | New RC? | Today's Lucid Tip: Dream Characters Suck | The RC List! | Failed Reality Checks | Funny dream cue / RC experience | The Automatization Technique | Reality Check Failure | Need a better reality check for Christ's sake!!!!!! | WHY?!?!?! Failed Reality Check! Threads about reality checks at The Lucidity Institute: Reality Testing September 2000 June 2001 August 2001 December 2001 Latest top&lt;br&gt; Techniques. When you read through these techniques, remember that different techniques work for different people. There is no “best technique” and most techniques could be used to have 2–5 lucid dreams every night! You could have several lucid dreams in a night, but you will not know it unless you remember them! However, you will probably want some advice as to which technique you should try first. Consider whether you want to use a method which starts from a dream or a method which starts from being awake.* If you master a technique which starts from being awake, you will eventually be able to have lucid dreams on demand when you sleep. For other techniques, you have to rely on your luck to give you lucid dreams after you have done your technique. Here are some advantages and disadvantages for specific techniques: WBTB. WBTB stands for “Wake-Back-To-Bed”. Wake yourself up after 4 to 6 hours of sleep, get out of bed and stay up for anywhere between a few minutes to an hour before going back to bed. It is preferable that you do something related to lucid dreaming during this time (such as reading about lucid dreaming), but it is not required. This is best combined with other techniques; many people have amazing results with a MILD/WBTB combination. The WBTB technique significantly increases your chances of a lucid dream, and using MILD (see below) in conjunction increases your success rate if you are planning to sleep an hour or more after your WBTB session. However, you might need plenty of sleep time and therefore you may only be able to use it on weekends. If you are sleeping too deeply to become lucid, then you can modify this technique. Try returning to sleep somewhere different than where you usually sleep, e.g. on a couch, a different bed, or even on the floor. If you are unable to do this, try changing the way you sleep, e.g. try sleeping with a lighter blanket or reversing the orientation of your body while in the bed (that is, swapping head and feet). Do this in order to teach your body that these different surroundings mean you want to have a more conscious sleep rather than a deeper sleep. In the beginning, different surroundings will also make you more alert, which can heighten your level of consciousness during sleep. I am sometimes awake for very short times, but cannot pull myself together enough to get up and out of bed. What can I do? Put a bright piece of paper on the wall or ceiling so that you will see it when you wake up. Other stimuli could be a hot water bottle, a light turned on under your bed, or an alarm clock. A good technique is to place an alarm clock out of arm's reach so that you are compelled to physically get up from bed and turn it off. If this is still insufficient to restore consciousness, try making a note of your intentions to remain awake and place the note on your alarm clock. After you get a lucid dream with this method, you'll find it easier and easier to get out of bed because you'll have more motivation. Threads about the WBTB technique at ld4all.com: The BIG WBTB topic Auto-suggestion. This technique describes how to use auto-suggestion to have lucid dreams. It can be especially effective for people who are highly susceptible to hypnosis or understand meditation, but for most people, MILD will probably be more effective. As you are falling asleep, suggest to yourself that you will have a lucid dream either that night or in the near future. You can use a mantra (such as “I will recognize that I'm dreaming”) if you want, but make sure you don't try too hard to get a lucid dream. Instead of putting intentional effort into the suggestion, try to genuinely expect to have a lucid dream. Let yourself think expectantly about the lucid dream you are about to have, but be patient if you don't get one right away. You may also use auto-suggestion to improve dream recall. Just use the technique as described above, but instead of suggesting that you'll have a lucid dream, suggest that you'll remember your dreams when you wake up. You could also use a mantra with this, such as “When I wake up, I will remember what I dreamt”. Just be careful not to put too much intentional effort into the mantra — try to genuinely expect to remember your dreams instead. MILD. MILD stands for “Mnemonic Induction of Lucid Dreams", or sometimes, “Mnemonically Induced Lucid Dream". The MILD technique was developed by Stephen LaBerge, and is described fully in his book "Exploring the World of Lucid Dreaming". With the MILD technique, as you are falling asleep, you concentrate on your intention to "remember" to recognize that you are dreaming. Repeat a short mantra in your head, such as “Next time I'm dreaming, I will remember I'm dreaming”. Think about what this means (i.e., that you want to remember that you are dreaming—in the same way you might go to a grocery store and suddenly remember that you need bread), and imagine that you are back in a dream you've had recently, but this time you recognize that you are dreaming. For example, if you recently dreamed of flying, imagine realizing that it is a dream because you are flying. Keep repeating and visualizing the mantra until you are sure that your intention is set in your mind or you fall asleep. If you stop repeating and visualizing the mantra, then still try to make sure the last thing in your mind before falling asleep is your intention to remember to recognize that you are dreaming. In general, the MILD technique can be practiced when you first go to bed at night, or after you have awakened from a dream during the night. If you practice the MILD technique after you have awakened from a dream, you should first run through the dream to ensure that you remember it. Some people find it helpful to jot down a few notes about their dream in their dream journal. Once you have committed the dream to memory, go back to sleep and follow the steps above. But this time, visualize the dream you just had. Move through the dream in your mind until you encounter a dream-sign you originally missed. Now, instead of missing the dream-sign, imagine yourself recognizing it and becoming “lucid”. Repeat this until you have fallen asleep. You "will" re-enter the dream, and you "will" recognize the dream-sign and finally, become lucid. Threads about the MILD technique at ld4all.com: The BIG MILD topic I II | MILD at Midnight? | MILD Mantras archive.org threads about the MILD technique from www.dreamviews.com (offline): Mnemonic Induction of Lucid Dreams (MILD) | Getting more help with MILD from your subconscious WILD. WILD stands for “Wake-Initiated Lucid Dream”, or “Wake-Initiation of Lucid Dreams” to refer to any technique that involves falling asleep consciously. These techniques are similar to self-hypnosis. Some people believe that WILDs are not actual dreams, but are instead astral projection. Various detailed resources are available under that moniker. For most people, they are far easier to induce in the early morning after waking up or in afternoon naps, as the sleep cycle will continue with a REM period. Once you are experienced with inducing WILDs, you can try to induce them at other times. For WILDs to occur, it is best for your body to be completely relaxed. When you go back to bed, lie down comfortably. Now tense and relax your body, starting from your shoulders and working downwards, then back up to the face. This (or similar relaxation, meditation, or trance techniques) should make your body feel slightly heavy and relaxed. There are many different ways to induce WILDs, but they all involve simultaneously attempting to keep the mind aware while attempting to have the body fall asleep. A few techniques are detailed below. If you pay attention to your physical body while using these techniques, then you will likely enter sleep paralysis (which usually happens after you are already asleep) without losing conscious awareness of your body. You will get a tingling and buzzing sensation (this might be unpleasant). These sensations might be so strong that you feel that you will die (e.g., you might feel a choking sensation), but don't worry, this is perfectly safe! In fact, this process happens to you every time you sleep, you are just not conscious during it. Sometimes you can simply wait until you fall asleep straight into a lucid dream. However, if you do not fall asleep, and you become completely paralyzed (with the exception of your eyes), do not try to move. Imagine your dream hand (or spirit hand if you prefer) going up and leaving your physical hand behind. Now you should have two separate bodies, a dream one and a real one. Control your dream body only — if you control your real one, you will wake up. Now you can try to roll out of bed into your dream world (alternatively, you can get up and walk through a mirror, or sink into your bed). There is a possibility that after waking up from a dream that you initiated using this technique, you may still be paralyzed. If this phenomenon occurs, it may be accompanied by hallucinations. For example, you may wake up from a lucid dream that you started using one of the WILD techniques, and you will still be paralyzed, this is a state of Sleep paralysis. If you have the experience then you have been given an example of what your mind can do, even without your direction. Before the phenomena was understood, as Sleep paralysis also occurs outside of Lucid Dream experiments it had been given a general negative outlook, while the experience is truly unique (partial awareness, unable to move, and multi-sense hallucinations) it is often linked to cultural fears, nightmares and other night terrors, Americans call it Old Hag Japanese people call it kanashibari, Korean people call it kawi nullida (being pressed by a pair of scissors), Turkish people call it "karabasan" and in old French it was called "Cauquemare", which was later transformed into the term "cauchemar". In any case as soon as you rationalize what is happening you should lose any fear by the temporary loss of muscular control and contextualize any hallucination. Chakra Technique. Using Chakra ("third eye") meditation is one way to achieve WILDs. Basically, one has to focus solely on his third eye and breathing to achieve a lucid dream. This is a technique over 5000 years old, taught to Parvati by Shiva. He is quoted as saying:With intangible breath in center of forehead, as this reaches heart at the moment of sleep, have direction over dreams and over death itself. Eyelid Patterns. This has been found to be very effective in many cases. However, it may lead to strange after effects, such as up to 15 dreams in one night, but otherwise, nothing harmful. This technique is as follows: When you go to sleep, ensure the room is completely dark. Then, with eyes closed, try to focus on the little dots that should appear to be moving on your eyelids. You will find that you can change their color at will. Continue focusing on these dots; make them dance around and form patterns and change colors, and eventually you should enter a Lucid Dream. This may take a little practice, but is usually very effective for summoning Lucid Dreams at will. Works better in conjunction with WBTB and other techniques mentioned above, but is still extraordinarily effective on its own. Hypnagogic Imagery. Stimulate your thinking patterns by constantly switching your attention. After doing this long enough, the images and sounds should begin to take a momentum on their own (this is called hypnagogic imagery), and may they get very strange and illogical. You should enter a dream at this point and quickly become lucid. Otherwise, you will eventually realize you have entered sleep paralysis consciously (see above). Because hypnagogic sleep paralysis involves full consciousness, dreaming can sometimes be frighteningly real. There is often a feeling of being flipped upside down, spun, or being tugged upon by an outside force. Hypnagogic hallucinations may also include strange auditory hallucinations, dark beings and flying. It is possible to observe waking reality while in a hypnagogic state, but this is limited to the sensations of your physical body. Most hypnagogic sleep paralysis states occur when sleeping face up. There is evidence that the tendency toward experiencing Hypnagogic sleep paralysis may be hereditary. Counting. Another technique is to count up to 100 in your head, optionally adding (for example) an “I'm dreaming” between each number. Alternatively, you can imagine stepping down stairs and reading each floor number, from 100 to 0. Try to make this image as vivid as possible — include not only what you see, but also what you hear, feel (touch the banister), and smell. At some point this image should continue into a dream or you will begin to get sleep paralysis as described above. It is easy to lose count, especially with counting up to 100 with an 'I'm dreaming' with each number. But stay focused: you are not going to sleep; your body is, and you must concentrate fully. Sound Technique. This method is for people, who can hear the "tinnitus". The idea is pretty much the same as the other WILD methods, which is to remain conscious while entering the dream state. In order to use this method, you must sleep in a perfectly quiet place. You need to hear the inner buzzing sound inside your ears. Lay your body down and relax as much as possible while trying to hear the sound. If you are too tired, you will fall asleep too fast and it will be difficult to remain conscious - in this case you should combine it with WBTB. By time you realize that the buzzing sound will increase in intensity. This might frighten newcomers, but be assured - nothing bad is going to happen. No, you will not be deaf when you wake up, it’s perfectly safe! *It is just an effect caused by your brain is trying to change mode, from listening to the ambient sound, to listening to the sound of dreamland, which is not real sound but just electrical charge inputed to the part of the brain to create a sensation of hearing. By that time, you will enter the hypnagogic state. All you need to do is concentrate, do not be afraid or think of anything, just be still, and in time your dream body will float, separating from your physical body, and there you go, you arrive in the dreamworld. Note: These sounds can be heard when you concentrate, even throughout the day - when you pay attention to them, but they shouldn't be heard much aloud if you are not in silence and concentrating on it - *as author said. The formula to hear it is quite simple: come into complete silence close your eyes and listen: is there only absolute silence or is there - something else? some...? do not create some sound just concentrate - :-) The best thing on this tech is - you can do it when you want. Slight Physical Discomfort. For the purpose of helping to retain your conscious awareness, slight physical discomfort is useful while performing any WILD technique. This prevents you from just drifting off to sleep. If you are lying down on your bed to do WILD and you are totally comfortable then your chances of going to sleep instead of remaining conscious are very high. The WILD technique relies on a form of deep trance induction, and many people who induce trances for other reasons rely on slight physical discomfort — for example the lotus position, or sitting in a hard-backed chair. Depending on your own preferences and your requirements of discomfort for success, you could choose from the following methods, arranged in ascending order of discomfort: Threads about the WILD technique at ld4all.com: The BIG WILD topic I II III IV V VI VII VIII IX X XI XII XIII XIV| Strange colored dots?!? ULTIMATE WILD method | Threads about the WILD technique at dreamviews.com: Wake-Initiated Lucid Dream (WILD) | Five Stages of WILD | WILD induction. help, only 3 hours til sleep time | Lotus-Flame WILD technique Incubating Dreams. To incubate a dream about a specific topic, you should first think of a phrase that summarizes that topic (e.g., “I want to go to Atlantis.”). It may help to write the phrase down. If there is something you want to do in the dream, think of a phrase to summarize that too (e.g., “I want to watch Atlantis sink into the ocean.”). If you want to become lucid in the dream, then you should probably write something like “When I dream of [the topic], I will remember that I'm dreaming.” beneath your topic phrase. Immediately go to sleep and focus on your topic phrase. Visualize yourself dreaming about the topic and (if you want to become lucid) realizing that you are dreaming. If there is something specific you want to do in the dream, visualize yourself doing it once you become lucid (not very likely to work if you don't become lucid in the dream). Think about your phrase and topic (and intention to become lucid) as you fall asleep. Make sure that the last thing in your mind before falling asleep is your intention to (lucidly) dream about the topic you want to dream about. You might want to wake yourself up when the dream starts to fade so that you remember more of the dream; you can do this by ignoring your perception of the dream environment — the opposite of dream stabilization techniques (just make sure you do a reality check when you wake up to make sure you are really awake). Chaining Dreams. Dream-chaining or “chaining dreams” is a method to re-enter your dream "after you've woken up". It can work for lucid and non-lucid dreams, but you will probably want to enter your dream lucidly. Once you wake up from a dream (if you don't think you were dreaming before you woke up, it may not work well) you should stay still and keep your eyes closed. Small movements are okay, but the less movement, sensory stimulation, and less time awake, the better. Ideally, it should feel less like you've woken up, and more like you've taken a 30 second break from dreaming. Once you are prepared to go back to sleep, close your eyes and either visualize yourself back in your dream, or use the “spinning technique” given in the next chapter to imagine yourself spinning back “into” your dream. Spinning is a little faster than visualization. Be sure to maintain the fact that you are dreaming (unless you don't want to be lucid), or you may lose your lucidity while falling asleep. Once in the dream, stimulate your senses as early as possible (see the next chapter). VILD. VILD stands for “Visual Induction of Lucid Dreams”, or sometimes, “Visually Incubated Lucid Dream”. This technique has been perfected by Peter Harrison, known as Pedro on the forums at ld4all.com. You may wish to read the main thread about the technique. The version described here has been adapted slightly. First, make sure you are relaxed. You can use the relaxing technique mentioned in the description of the WILD technique. You can also imagine your brain emptying out and becoming sleepier. If you have a hard time falling asleep quickly, it should help to read a book (preferably about lucid dreaming) for a while before you go to sleep, until you feel very sleepy. Now, you need to visualize a dream which you had prepared earlier. Here is an example of a prepared dream: Make sure you know exactly what the dream would be like, such as which friend, the exact words they say, and which reality checks you do. Reality checks that require no props, such as books or clocks, are recommended. Visualize this dream slowly three times, to make sure that you know every detail. Then, start going full-on and visualize the dream over and over. You should visualize the dream as though you are looking through your own eyes, not from a third-person perspective. If you find your thoughts drifting, ignore them and continue to visualize the dream continuously. You will need patience for this — don't just give up if you think it won't work. When you actually dream this, you will not notice the difference — until you do your reality checks! Continue with the dream as you incubated it (e.g., remember to thank your friend!) before continuing through the door. I tried to visualize the dream until I fell asleep, but I just stayed awake. What went wrong? If visualizing keeps you awake, the VILD technique is not the technique for you! You should use a different technique. Threads about VILD at ld4all.com: I can LD at will!!!! I II III | VILD...Visually Incubated Lucid Dream Topics about VILD at The Lucidity Institute: The VILD Technique There is an appendix on VILD. LILD. LILD stands for “Lucid Induction of Lucid Dreams”, or sometimes, “Lucidly Induced Lucid Dream”. To use this technique, you need to have a lucid dream in the first place, but it can help you to get more later. The idea is to do something in your dream that will help you to become lucid the next time you are dreaming. For example, you could ask a dream character for help — ask them to meet you the next night and tell you that you are dreaming. If it works out the way it should, then the next time you are dreaming, the dream character will walk up to you and tell you that you are dreaming, and so you'll (hopefully) become lucid. There are many variations on this technique; you could set up signs in your dreamworld that remind you to do a reality check or eat lucid pills instead! This technique is not likely to be very effective, but it "can" work; it relies on the chance that you'll subconsciously induce the reminder (i.e., the dream character or sign or whatever you used) during some later dream, and become lucid because of it. Note that LILD is best used in conjunction with dream-signs and auto-suggested non-lucid dreams. The basic idea as explained above is to have something in your dream that triggers the transition from normal dream state to lucid dreaming. To simply tell a character to tell you that you are dreaming the next time you fall asleep is usually not enough. There is no guarantee that you will dream about that character and there is no guarantee that your subconscious will believe the character enough to make you snap into lucidity (make you realize that you are in fact dreaming). Now as this technique suggests, you must have some previous alternate means of having a lucid dream. Whatever technique you employ to get into this initial lucid dream state is not really important, but you should try to remember to use this technique (LILD) once you do get into a lucid dream state. Thinking of this before falling asleep (MILD) sometimes helps and usually takes many lucid dreams before finally remembering. Once you are in a lucid dream, make up a dream-sign. It can be anything. It can be an object. It can be food or a drink (that doesn't taste like anything). It is usually best to pick something that isn't quite right. Something that on the surface would appear normal in the real world, but that upon closer inspection is not quite right. Food or drinks are good as they can have no taste or not be refreshing in a dream. But try and pick something that you dream about a lot so that there is a better chance of you dreaming about this dream-sign later on. Now pick something else that only appears or happens in your lucid dream. It can be anything. If there is nothing in your current lucid dream, create something really strange. Something that could never be confused with the real world. Now mentally associate the dream-sign (food) with this unusual item or event that could never happen in the real world. But at the same time, this unusual item or event should equate to "lucid dreaming". When you see the unusual item, it should only make you think of when you have a lucid dream as this should be the only time you encountered it. So we have a 3 item associative link. Do all of the above while in a lucid dream. The next time you dream about your dream-sign, your subconscious will think of the unusual item or event. The unusual item or event will make you think of lucid dreaming. The two combined impossibilities (1. dream-sign that cannot exist in the real world 2. item or event that only appears in lucid dreams) will make your unconscious try to make a decision on all this. This will make your conscious mind come to the surface and hopefully you will come to the conclusion that you are dreaming. Many times, you will not want to deal with it because you are too tired (that is why you are sleeping, no?) and fall back into a normal dream state. This is why it can take a few tries. Eventually, your subconscious will start putting clear signs in your dreams like billboards that spell out "YOU ARE DREAMING". But once it triggers, it is quite the realization that an instant before, you had no real control over your actions and now you can do whatever you want. Another note... if it failed, you will usually know why. So next time, you can choose another dream-sign or slightly different technique or something more shocking. Once you get this working once, it is relatively easy to use over and over as the hard part just described is over with. Sometimes disassociative techniques are needed if used too much. To sum up, this technique is a way to force a reality check while in a normal dream state where your subconscious has no choice but to come to the conclusion that you are in fact dreaming. Once your mind knows that you are dreaming, there will be no other conclusion than your conscious mind taking over. And this is what lucid dreaming is all about. CAT. For detailed information on the Cycle Adjustment Technique, see the appendix on CAT. Topics about CAT at The Lucidity Institute: **CAT method** New Lucid dream induction technique There is an appendix on CAT. Tibetan methods. Tibetan Buddhists practice what is known as Tibetan dream yoga. Probably the most time consuming way of inducing lucid dreams, it is also, according to the practitioners, the most rewarding. The basic practice is "awareness". Awareness should be practiced while sleeping just as well as while being awake. Meditating on the question “who is aware?” might catapult you into a higher degree of awareness according to Buddhist beliefs. Keeping this level of awareness is another matter. The Tibetans have developed many yogic exercises and disciplines to be practiced. Maybe the most interesting difference between Tibetan dream yoga and modern western methods of lucid dream induction is the Tibetan claim of the possibility to be aware during deep sleep, not only in the REM periods of sleep. For the reader who is interested in these methods a good start is to begin to regard all experience as a dream. After all, from the countless multitude of matter and radiation reaching our senses the nervous system tunes in only to a small fraction of this chaos. For members of the phalanx that believes we, more or less, create our own reality in the above sense this practice should feel natural. In general though, it is recommended to gain instruction from a teacher in the flesh rather than from books (like this one!). The Tibetan Methods are not "time consuming" if the goal is to go into deeper than Lucid states. From a Raja Yoga stand point, from a Daoist standpoint, and from a Tibetan Yoga standpoint the goal is not to "play" in Lucid Dreams, but to dispel the delusional nature of what we call "reality." Also, to heal and to develop super-learning, and to increase the "energy" you find most attractive: innovative, creative, joyous, blissful, love, power, wisdom, mirth. Are you Ambitious-Worldly or Ambitious-Spiritual? All these factors play into it. These techniques are very easy to learn and are rooted in a deeper science than the WILD, VILD, LILD -- one that does not use auto-suggestion or forced recall. In conclusion, the Tibetan methods are advanced yogic techniques that only an accomplished practitioner can fully teach. If one does not have access to a Tibetan dream yogi, one should probably concentrate on techniques that are more mundane from a Western point of view. Other techniques. Many of these are combinations of other techniques with some addition or modification. top&lt;br&gt; Other methods. Food and drink. There are various foods and drinks that you can consume which seem to have some effect on sleeping and dreaming. Note that for most of these there is no explanation or scientific study of how they work, and some may simply be placebos. Don't go overboard with the consumption of any of these, as overdosing could have nasty effects (well, milk should be safe unless you are allergic). Don't experiment without accumulating enough knowledge first. The authors in no way encourage the use of legal or illegal drugs. Some people who practice Lucid dream (LD) or Out of body experience (OBE) use Galantamine to increase their odds to achieve LD or OBE. By taking small amount of Galantamine (around 4 to 8 mg) after 5 to 6 hours of deep sleep and practice the induction technique such as meditation, MILD or WILD many people report more success with Galantamine. There also report that taking Galantamine without proper induction technique will not lead to LD or OBE but will result in only a vivid dream instead. It should also be noted that due to a long half life Galantamine will stay in the body for a period of up to and over 48 hours, as such it is advisable to space out the use of Galantamine over a period of three days so that the body does not build a resistance to the drug ruining its effectiveness. Some people report mixing Galantamine with other Nootropic can enhance the degree of lucidity but this is still controversial since some mixtures may work for some people, but lead to failure for others. Plants: For an in-depth guide to using supplements for lucid dreaming, see the book Advanced Lucid Dreaming - The Power of Supplements. Drugs. Dissociatives and hallucinogens can be used to create a (more or less) lucid dream-like state, though whether or not these help with lucid dreaming is debatable. The authors do not recommend use of these substances for induction of lucid dreams, nor do they condone the breaking of any applicable laws. Some dissociatives and hallucinogens are: For more info, see Erowid Vaults Gadgets. There are various gadgets you can use to become lucid easily. They generally detect when you are in the REM state and then provide a light and/or sound signal. This signal is supposed to be adjusted so that it doesn't wake you up but does enter your dream. The signal is then recognized as showing that you are dreaming, and you become lucid. The most well-known device is the NovaDreamer from the Lucidity Institute. However, this product is no longer produced. Be sure to check for recommendations for devices from lucid dreaming forums. A similar device is the DreamMaker. The DreamMaker works very similarly to the NovaDreamer but without the Dream Alarm feature, which worked to wake the dreamer in the middle of the REM state. This device comes with a mask, a circuit board with adjustable controls, the batteries needed to operate it, a short owner's manual, a lucid dreaming workbook, and the Stephen LaBerge book "Exploring the World of Lucid Dreaming". The circuit board is supplied completely ready to use, but you have to insert the batteries and put the circuit board into the mask yourself. An alternative is the Kvasar. The Kvasar costs about $20 in raw materials, but needs to be constructed by somebody skilled in electronics as it is not sold commercially. It can also be hard to operate. Another do-it-yourself alternative to commercial dreaming masks is Nate True's Lucid Dream Mask, which does not bother with difficult-to-calibrate sensors and just uses a timer for flashing lights, and has (ostensibly) competitive results with all of the former gadgets. The owners of Wellness Tools, who makes the DreamMaker, and Kvasar have not had friendly relations; see DreamViews.com. Software. There are many programs for your computer that can assist with lucid dreaming. These can give out verbal cues while you sleep, or assist in doing your reality checks: There is also a list of programs available at LD4all.com, under the “How” section. Hypnosis. The techniques commonly referred to as hypnosis is based in a set of hypnotic suggestions are commonly composed of a series of instructions and ideas that may be delivered by a hypnotist in the presence of the subject, or they may be self-administered ("self-suggestion" or "autosuggestion"), even in a conscious state. Hypnosis is not a science. In fact, it has in almost every aspect eluded scientific analysis, as it is extremely hard to generalize (each individual responds differently and at different levels) and the methodology is so diverse and based on yet to be completely understood mental/biological phenomena, most related with faith or the placebo effect. For the phenomena to work, whom in aggregate we define as hypnosis, the mind has to be able to turn the suggestions into reality. All hypnosis is ultimately self hypnosis. If you, for instance, take into consideration the problems that faith beliefs have caused to the human race, or even the problems in the field of psychology, you can appreciate the problem of scientifically studying hypnotic phenomena, as they are extremely open to individual interpretation and to one's ability to be open to suggestibility (self-induced or otherwise). The Wikibook's work on Hypnosis will covers the subject in greater detail. It will also offer support material, or reference it, such as methods to induce, or reinforce specific mental states (or states of mind) and any attitudes and beliefs that impact the phenomena. Including information on techniques or procedures for hypnotic induction and hypnotic suggestion. References. top 

Brake and derailleur cables are extremely difficult to cut with a standard pair of wire cutters (sometimes called side or diagonal cutters); when it is finally accomplished, this process usually leaves a frayed mess of the inner cable, which is impossible to slip through the cable housing or the hole in a fastening bolt. In addition, the jaws of these cutters are made for cutting copper electrical wire, and will become notched when cutting steel bicycle cables. Cable cutters specifically designed for steel cable, on bicycles or elsewhere, are available at prices low enough to eliminate this nuisance. They are generally similar to regular cutters or scissors, but have extra-hard jaws, each of which carries a v-shaped notch. When cutting a cable, the notches overlap, trapping the cable in a diamond shaped hole with sharp edges, which becomes smaller as the handle is compressed, thereby cutting the cable easily without fraying and with no damage to the tool. They are also useful for cutting the cable housing, although they often close the inner sleeve, and it needs to be re-opened with a piece of cable or other means. Sometimes their use leaves a small burr on the housing which may need to be filed off. 

 Welcome to Cryptography, the study of obfuscating data to unintended recipients. Part I: Introducing Cryptography Part II: Designing Cryptosystems Part III: Cryptanalysis Part IV: Using Cryptosystems Part V: Cryptography and Society Part VI: Miscellaneous Pages to be merged into the text. Cryptography/Prime Curve/Affine Coordinates Cryptography/Prime Curve/Chudnovsky Coordinates Cryptography/Prime Curve/Jacobian Coordinates Cryptography/Prime Curve/Standard Projective Coordinates Cryptography/Notes 



Canada's Obligation. By the 1860s, Great Britain became really concerned that Canada was more of a shining star than an asset to the Empire. British politicians did not like the expense of providing for the defence of the North American colonies from real or perceived threats from the Americans. The American Civil War only heightened fears in Canada of a potential invasion. Americans, on the other hand, always seemed to have a group of people willing to go to war with Britain and invade British North America, often in high places (for example, Lincoln's Secretary of State, Seward). Britain's trade with the United States at this time had simply become too valuable to jeopardize by maintaining a colonial system in British North America, so prominent British politicians began the process of pushing Canada "out of the nest," so to speak. However, the process was not simple. British North America consisted of no less than seven colonies, and the vast territory of Rupert's Land owned and run by Hudson's Bay Company. Uniting this area would prove to be difficult. --- The following passage is inaccurate -- Confederation was achieved for many reasons. Canada West and Canada East felt that uniting the colonies would help make all of the colonies stronger, more economically stable, and would make the government system more fair. The maritimes, who were having a boom period, were disinclined to join confederation. They wanted to join together the maritime colonies and have a separate country for only them, however, the Americans had just finished the Civil War, and had about as many soldiers as all the colonies that were thinking about confederation did combined, and posed a real threat to the colonies. The maritimes thought about it, and with the help of a conference with John A. Macdonald decided to join confederation. The only maritime colonies that didn't join were Newfoundland and Prince Edward Island, which would join later on. At the Quebec conference, representatives from all the participating colonies came together and constructed the 72 Resolutions which outlined all of the laws that the new country Canada would have. 

Dream stabilization. Once you are able to dream lucidly, you may find that it is difficult to stay in the dream; for example, you may wake instantly or the dream may start “fading” which is characterized by loss or degradation of any of the senses, especially vision. Alternatively, a new lucid dreamer could easily forget that they are in a dream, as a result of the shock of the sensation. Don't worry if you wake immediately after becoming lucid. As you gain more experience of becoming lucid, it will come as less of a shock and you’ll be less likely to wake up. Make sure you do a reality check to be sure you’re not still dreaming. As you gain more experience, you will have an easier time identifying when and remembering that you are dreaming. You can avoid more gradual fadings by stimulating your senses. This means listening for sounds, feeling around with your hands, and paying attention to what you see and smell. The idea here is to load your senses with stimulation from the dream so that your senses cannot shift to the real world. If you close your eyes, you are removing a great deal of sensory information and might wake up. Staring at a single point can cause effectively the same problem if you stop seeing everything else in your peripheral vision, or don't see enough movement. If you hear something loud in real life and are hearing nothing in the dream, your senses may shift to the real world, causing you to wake up. Ideally you should be able to use the techniques below to stabilize your dream "before" it starts to fade (or “black out”). As always, prevention is better than treatment - and the more stable and vivid your dreams are, the more enjoyable they will be. However, if that doesn't work you may be able to use stabilization techniques to stop the fading; the spinning technique is probably the most effective in this case. If you still can’t stabilize your dream, you may decide to try and wake up with the aim of remembering your dream as accurately as possible while its still fresh in your mind. Hand Touching. Rub your hands together and concentrate on the rubbing. You should feel the friction and the heat of your hands. If you can concentrate on the feelings that this action generates, your dream is likely to stabilize and become more vivid and detailed. You can also keep one hand on your arm while exploring the dream for a constant sense of stimulation. This technique is most effective when used in conjunction with the “Slowing it down” technique, by staring at your hands while rubbing them together. Spinning. You spin around in your dream much as you would if you suddenly want to feel dizzy in real life. The sensation of movement is the key here to stabilizing the dream. Many people report success with this technique, but it also tends to cause a complete change of your dream scene (see Changing the dream environment below). If the dream scene disappears completely (e.g., becomes black), it is necessary to visualize the dreamscape to return to the dream. Slowing it down. Some people like to stabilize the dream by “stopping to smell the roses” and slowly staring at a dream object until it becomes clear. The dreamer would then look around elsewhere, noticing how detailed everything is, thereby stimulating the visual portion of the dream. "However, others find this can cause their lucid dream to end." If you focus on one object for too long to the exclusion of everything else, you will likely wake up or lose the dream. It works best to pay attention to everything in your vision, including your peripheral vision, not just the center of the object you're staring at. If staring at a single object doesn't work for you, try to let your eyes wander around instead. Touching your dream. If you feel that your dream is too abstract and fear that it might be fading, you can prevent this by grabbing hold of a solid object in your dream and focus on how real the sensation is. A good tip is to find something you know is stuck, for instance a table nailed to the ground, and pull it with all your muscular power (no "supernatural" powers!), and you should feel how solid it is. The idea is that you convince yourself that the dream is solid and real — through tactile stimulation — and nothing abstract. Regaining waking memory or skills. This is also likely to enhance your degree of lucidity. Try to remember facts from your waking life, such as your phone number, address, etc., or do some simple math. Or, start reciting the lyrics to your favorite song. Or perhaps try some sports practice you know well — this all depends on which senses / methods of thought process you tend to rely on most in your waking life. Dream Guide. Summon a DC (Dream character) that you expect to Guide you through the dream. Summoning is easier if you always expect this Guide to be there in a dream. The Guide is useful for escaping scary or unpredicted environments ("whooshing" you away, killing or unsummoning the threat, reminding you that you are in control.). Imagine this Guide as powerful, knowledgeable, and protective over you. Guides can also walk you through summoning, flying, shapeshifting, etc. and be useful for learning and controlling your dreams. False awakening. A couple of the users on the ld4all.com forums have had success with creating a false awakening to stabilize a dream. If the above techniques are failing and you find your dream still fading, "and you really want to continue your lucid dream", do the following: You will either have a false awakening, reality check, and then end up with an even more vivid lucid dream, or will really wake up, perform a reality check, and realize that you just woke up (bad luck!). The most important part of this is the reality check. This is what will continue your lucid dream. You should be performing reality checks when you wake up. If you plan to induce false awakenings in order to stabilize a dream, the reality check that you perform as you wake up is as important as the one that got you lucid, if not more. Perform every check in the book until you are positively, absolutely, and completely sure that you aren’t dreaming. A series of 10 reality checks is more likely to produce dream results in a dream, especially if you are expecting dream results. This technique is for those who are desperate! "If you have had a good experience with this technique, please go to the talk page and post your experiences, as there have not been many anecdotes of it working yet." top Recovering from lost visuals. There are a few things you can try to do if you lose your vision. Most of these are less likely to help prolong your dream than the above techniques. You can also try these if you have just woken up and are lying in your bed. You may be able to return to your dream. Autosuggestion. You can repeat over and over a phrase similar to “I can see my dream,” or otherwise enforce in your mind that you can see a dreamscape. (See Autosuggestion) Visualising. You can visualise the scene as it would be if you could see it. You could take this as an opportunity to change the dreamscape by visualising a different environment from the previous one in the dream. This can be made easier by spinning as you visualize. (See Changing the dream environment below) Altering the dream. Changing the dream environment. You can change the dreamscape by simply visualizing a different environment. Stephen LaBerge, author of "Exploring the World of Lucid Dreaming", suggests closing your eyes, spinning around, and visualizing a new location. Alan Worsley, a famous lucid dreamer, describes another technique of summoning a television and remote to switched dreamscapes. By simply “changing the channel” on the remote, imagine your surroundings switching to your desired location, as though you were switching through various television programs. In both of these methods, details are key. The more details you provide for your next dreamscape, the easier it will be to get there. For example, if you wanted go to the Superbowl in your dreams, you might simply state: But this statement could be improved by adding details: Summoning objects into your dream. At some point in your lucid dreaming experience, you'll probably want to handle various objects, or talk with certain people. Both of these needs can be fulfilled by using your mind, and the power of suggestion. There is no set way to "make" things in the dream world; in fact, many lucid dreamers devise their own methods through experimentation. However, here are some of the more common "summoning" techniques: Remember, in the dream world, your expectations shape your surroundings. If you think a big, scary monster is going to step out of the shadows and attack you...well, a big, scary monster probably "is" going to step out of the shadows and attack you. So, don't be a victim; take control of your thoughts and use them to your advantage. top What you can do. This final section should see you off with a few ideas of what to do in a dream. It is advised to have a "clear purpose" for your lucid dreams whenever you go to sleep. In other words, every night you consider what you want to do when you have a lucid dream, and select one thing, or perhaps two or three if you are skilled. Avoid this: You will either end up doing none of these things in your dream or getting overexcited and waking up. Now that that’s clear, here’s a list of possible things you could do, ordered in difficulty. Remember that you might find some things unusually hard (or easy) compared to most lucid dreamers, this is perfectly normal! This is a very rough guide — if you’ve managed something in the Easy section, don't be scared to try for something from the Medium section. Hard. top Conclusion. With all the techniques in this book, you may feel overwhelmed and uncertain of what to do next. Don't worry — just choose a few techniques to “map your way to lucidity”, decide on a few things you will want to do from this page, and start! If you are still unsure of what to do, don’t worry — you might happen to have a lucid dream tonight! If you are beginning to feel a compulsive thirst for "more" information about dreams, head over to the Further Reading section for the sites to satisfy your cravings. Remember to come back occasionally and help make the wikibook grow! top 

It seems likely that in the distant past Vietnamese shared more characteristics common to other languages in the Austroasiatic family, such as an inflectional morphology and a richer set of w:consonant clusters, which have subsequently disappeared from the language. However, Vietnamese appears to have been heavily influenced by its location in the Southeast Asian sprachbund—with the result that it has acquired or converged toward characteristics such as isolating morphology and w:tonogenesis. These characteristics, which may or may not have been part of proto-Austroasiatic, nonetheless have become part of many of the philologically unrelated languages of Southeast Asia—for example, Thai (one of the w:Tai-Kadai languages), Tsat (a member of the Malayo-Polynesian language group within Austronesian), and Vietnamese each developed tones as a phonemic feature, although their respective ancestral languages were not originally tonal. The Vietnamese language has similarities with Cantonese in regard to the specific intonations and unreleased plosive consonant endings, a legacy of archaic Chinese that can also be found in Korean. The ancestor of the Vietnamese language was originally based in the area of the Red River in what is now northern Vietnam, and during the subsequent expansion of the Vietnamese language and people into what is now central and southern Vietnam (through conquest of the ancient nation of Champa and the Khmer of the Mekong delta in the vicinity of present-day Ho Chi Minh City), Vietnamese was linguistically influenced primarily by Indic and Malayo-Polynesian languages at first, until Chinese came to politically predominate the Vietnamese area toward the middle of the first millennium AD. With the rise of Chinese political dominance came a radical importation of Chinese vocabulary and grammatical influence. As Chinese was, for a prolonged period, the only medium of literature and government, as well as the primary language of the ruling class in Vietnam, much of the lexicon of Vietnamese in all realms consists of Hán Việt (Sino-Vietnamese) words. In fact, as the vernacular language of Vietnam gradually grew in prestige toward the beginning of the second millennium, the Vietnamese language was written using Chinese characters (see Chữ nôm) adapted to write Vietnamese, in a similar pattern as used in Japan (see "kanji"), Korea and other countries in the Chinese cultural sphere. The Nôm writing reached its zenith in the 18th century when many Vietnamese writers and poets composed their works in Chữ Nôm, most notably Nguyễn Du and Hồ Xuân Hương (dubbed "the Queen of Nôm poetry"). As contact with the West grew, the Quốc Ngữ system of Romanized writing was developed in the 17th century by Portuguese and other Europeans involved in proselytizing and trade in Vietnam. When France invaded Vietnam in the late 19th century, French gradually replaced Chinese as the official language in education and government. Vietnamese adopted many French terms, such as đầm (dame, from "madame"), ga (train station, from "gare"), and va-li (valise). In addition, many Sino-Vietnamese terms were devised for Western ideas imported through the French. However, the Romanized script did not come to predominate until the beginning of the 20th century, when education became widespread and a simpler writing system was found more expedient for teaching and communication with the general population. 

Cryptography is the study of information hiding and verification. It includes the protocols, algorithms and strategies to securely and consistently prevent or delay unauthorized access to sensitive information and enable verifiability of every component in a communication. Cryptography is derived from the Greek words: kryptós, "hidden", and gráphein, "to write" - or "hidden writing". People who study and develop cryptography are called cryptographers. The study of how to "circumvent" the use of cryptography for unintended recipients is called cryptanalysis, or codebreaking. Cryptography and cryptanalysis are sometimes grouped together under the umbrella term cryptology, encompassing the entire subject. In practice, "cryptography" is also often used to refer to the field as a whole, especially as an applied science. At the dawn of the 21 century in an ever more interconnected and technological world cryptography started to be ubiquitous as well as the reliance on the benefits it brings, especially the increased security and verifiability. Cryptography is an interdisciplinary subject, drawing from several fields. Before the time of computers, it was closely related to linguistics. Nowadays the emphasis has shifted, and cryptography makes extensive use of technical areas of mathematics, especially those areas collectively known as discrete mathematics. This includes topics from number theory, information theory, computational complexity, statistics and combinatorics. It is also a branch of engineering, but an unusual one as it must deal with active, intelligent and malevolent opposition. An example of the sub-fields of cryptography is steganography — the study of hiding the very "existence" of a message, and not necessarily the "contents" of the message itself (for example, microdots, or invisible ink) — and traffic analysis, which is the analysis of patterns of communication in order to learn secret information. When information is transformed from a "useful form" of understanding to an "opaque form" of understanding, this is called encryption. When the information is reverted back into a useful form, it is called decryption. Intended recipients or authorized use of the information is determined by whether the user has a certain piece of secret knowledge. "Only" users with the secret knowledge can transform the opaque information back into its useful form. The secret knowledge is commonly called the key, though the secret knowledge may include the "entire process" or algorithm that is used in the encryption/decryption. The information in its useful form is called plaintext (or cleartext); in its encrypted form it is called ciphertext. The algorithm used for encryption and decryption is called a cipher (or cypher). Common goals in cryptography. In essence, cryptography concerns four main goals. They are: Not all cryptographic systems achieve all of the above goals. Some applications of cryptography have "different" goals; for example some situations require repudiation where a participant can plausibly deny that they are a sender or receiver of a message, or extend this goals to include variations like: Common forms of cryptography. Cryptography involves all legitimate users of information having the keys required to access that information. Other: Poorly designed, or poorly implemented, crypto systems achieve them only by accident or bluff or lack of interest on the part of the opposition. Users can, and regularly do, find weaknesses in even well-designed cryptographic schemes from those of high reputation. Even with well designed, well implemented, and properly used crypto systems, some goals aren't practical (or desirable) in some contexts. For example, the sender of the message may wish to be anonymous, and would therefore deliberately choose not to bother with non-repudiation. Alternatively, the system may be intended for an environment with limited computing resources, or message confidentiality might not be an issue. In classical cryptography, messages are typically enciphered and transmitted from one person or group to some other person or group. In modern cryptography, there are many possible options for "sender" or "recipient". Some examples, for real crypto systems in the modern world, include: When confusion on these points is present (e.g., at the design stage, during implementation, by a user after installation, or ...), failures in reaching each of the stated goals can occur quite easily—often without notice to any human involved, and even given a perfect cryptosystem. Such failures are most often due to extra-cryptographic issues; each such failure demonstrates that good algorithms, good protocols, good system design, and good implementation do not alone, nor even in combination, provide 'security'. Instead, careful thought is required regarding the entire crypto system design and its use in actual production by real people on actual equipment running 'production' system software (e.g., operating systems) -- too often, this is absent or insufficient in practice with real-world crypto systems. Although cryptography has a long and complex history, it wasn't until the 19th century that it developed anything more than ad hoc approaches to either encryption or cryptanalysis (the science of finding weaknesses in crypto systems). Examples of the latter include Charles Babbage's Crimean War era work on mathematical cryptanalysis of polyalphabetic ciphers, repeated publicly rather later by the Prussian Kasiski. During this time, there was little theoretical foundation for cryptography; rather, understanding of cryptography generally consisted of hard-won fragments of knowledge and rules of thumb; see, for example, Auguste Kerckhoffs' crypto writings in the latter 19th century. An increasingly mathematical trend accelerated up to World War II (notably in William F. Friedman's application of statistical techniques to cryptography and in Marian Rejewski's initial break into the German Army's version of the Enigma system). Both cryptography and cryptanalysis have become far more mathematical since WWII. Even then, it has taken the wide availability of computers, and the Internet as a communications medium, to bring effective cryptography into common use by anyone other than national governments or similarly large enterprise. External Links 

Classical Cryptography. The earliest known use of cryptography is found in non-standard hieroglyphs carved into monuments from Egypt's Old Kingdom (ca 4500 years ago). These are not thought to be serious attempts at secret communications, however, but rather to have been attempts at mystery, intrigue, or even amusement for literate onlookers. These are examples of still another use of cryptography, or of something that looks (impressively if misleadingly) like it. Later, Hebrew scholars made use of simple Substitution ciphers (such as the Atbash cipher) beginning perhaps around 500 to 600 BCE. Cryptography has a long tradition in religious writing likely to offend the dominant culture or political authorities. The Greeks of Classical times are said to have known of ciphers (e.g., the scytale transposition cypher claimed to have been used by the Spartan military). Herodutus tells us of secret messages physically concealed beneath wax on wooden tablets or as a tattoo on a slave's head concealed by regrown hair (these are not properly examples of cryptography per se; see secret writing). The Romans certainly did (e.g., the Caesar cipher and its variations). There is ancient mention of a book about Roman military cryptography (especially Julius Caesar's); it has been, unfortunately, lost. In India, cryptography was apparently well known. It is recommended in the Kama Sutra as a technique by which lovers can communicate without being discovered. This may imply that cryptanalytic techniques were less than well developed in India ca 500 CE. Cryptography became (secretly) important still later as a consequence of political competition and religious analysis. For instance, in Europe during and after the Renaissance, citizens of the various Italian states, including the Papacy, were responsible for substantial improvements in cryptographic practice (e.g., polyalphabetic ciphers invented by Leon Alberti ca 1465). And in the Arab world, religiously motivated textual analysis of the Koran led to the invention of the frequency analysis technique for breaking monoalphabetic substitution cyphers sometime around 1000 CE. Cryptography, cryptanalysis, and secret agent betrayal featured in the Babington plot during the reign of Queen Elizabeth I which led to the execution of Mary, Queen of Scots. And an encrypted message from the time of the Man in the Iron Mask (decrypted around 1900 by Étienne Bazeries) has shed some, regrettably non-definitive, light on the identity of that legendary, and unfortunate, prisoner. Cryptography, and its misuse, was involved in the plotting which led to the execution of Mata Hari and even more reprehensibly, if possible, in the travesty which led to Dreyfus' conviction and imprisonment, both in the early 20th century. Fortunately, cryptographers were also involved in setting Dreyfus free; Mata Hari, in contrast, was shot. Mathematical cryptography leapt ahead (also secretly) after World War I. Marian Rejewski, in Poland, attacked and 'broke' the early German Army Enigma system (an electromechanical rotor cypher machine) using theoretical mathematics in 1932. The break continued up to '39, when changes in the way the German Army's Enigma machines were used required more resources than the Poles could deploy. His work was extended by Alan Turing, Gordon Welchman, and others at Bletchley Park beginning in 1939, leading to sustained breaks into several other of the Enigma variants and the assorted networks for which they were used. US Navy cryptographers (with cooperation from British and Dutch cryptographers after 1940) broke into several Japanese Navy crypto systems. The break into one of them famously led to the US victory in the Battle of Midway. A US Army group, the SIS, managed to break the highest security Japanese diplomatic cipher system (an electromechanical 'stepping switch' machine called Purple by the Americans) even before WWII began. The Americans referred to the intelligence resulting from cryptanalysis, perhaps especially that from the Purple machine, as 'Magic'. The British eventually settled on 'Ultra' for intelligence resulting from cryptanalysis, particularly that from message traffic enciphered by the various Enigmas. An earlier British term for Ultra had been 'Boniface'. World War II Cryptography. By World War II mechanical and electromechanical cryptographic cipher machines were in wide use, but they were impractical manual systems. Great advances were made in both practical and mathematical cryptography in this period, all in secrecy. Information about this period has begun to be declassified in recent years as the official 50-year (British) secrecy period has come to an end, as the relevant US archives have slowly opened, and as assorted memoirs and articles have been published. The Germans made heavy use (in several variants) of an electromechanical rotor based cypher system known as Enigma. The German military also deployed several mechanical attempts at a one-time pad. Bletchley Park called them the Fish cyphers, and Max Newman and colleagues designed and deployed the world's first programmable digital electronic computer, the Colossus, to help with their cryptanalysis. The German Foreign Office began to use the one-time pad in 1919; some of this traffic was read in WWII partly as the result of recovery of some key material in South America that was insufficiently carefully discarded by a German courier. The Japanese Foreign Office used a locally developed electrical stepping switch based system (called Purple by the US), and also used several similar machines for attaches in some Japanese embassies. One of these was called the 'M-machine' by the US, another was referred to as 'Red'. All were broken, to one degree or another by the Allies. Other cipher machines used in WWII included the British Typex and the American SIGABA; both were electromechanical rotor designs similar in spirit to the Enigma. Modern Cryptography. The era of modern cryptography really begins with Claude Shannon, arguably the father of mathematical cryptography. In 1949 he published the paper Communication Theory of Secrecy Systems in the Bell System Technical Journal, and a little later the book "Mathematical Theory of Communication" with Warren Weaver. These, in addition to his other works on information and communication theory established a solid theoretical basis for cryptography and for cryptanalysis. And with that, cryptography more or less disappeared into secret government communications organizations such as the NSA. Very little work was again made public until the mid '70s, when everything changed. 1969 saw two major public (i.e., non-secret) advances. First was the DES (Data Encryption Standard) submitted by IBM, at the invitation of the National Bureau of Standards (now NIST), in an effort to develop secure electronic communication facilities for businesses such as banks and other large financial organizations. After 'advice' and modification by the NSA, it was adopted and published as a FIPS Publication (Federal Information Processing Standard) in 1977 (currently at FIPS 46-3). It has been made effectively obsolete by the adoption in 2001 of the Advanced Encryption Standard, also a NIST competition, as FIPS 197. DES was the first publicly accessible cypher algorithm to be 'blessed' by a national crypto agency such as NSA. The release of its design details by NBS stimulated an explosion of public and academic interest in cryptography. DES, and more secure variants of it (such as 3DES or TDES; see FIPS 46-3), are still used today, although DES was officially supplanted by AES (Advanced Encryption Standard) in 2001 when NIST announced the selection of Rijndael, by two Belgian cryptographers. DES remains in wide use nonetheless, having been incorporated into many national and organizational standards. However, its 56-bit key-size has been shown to be insufficient to guard against brute-force attacks (one such attack, undertaken by cyber civil-rights group The Electronic Frontier Foundation, succeeded in 56 hours—the story is in "Cracking DES", published by O'Reilly and Associates). As a result, use of straight DES encryption is now without doubt insecure for use in new crypto system designs, and messages protected by older crypto systems using DES should also be regarded as insecure. The DES key size (56-bits) was thought to be too small by some even in 1976, perhaps most publicly Whitfield Diffie. There was suspicion that government organizations even then had sufficient computing power to break DES messages and that there may be a back door due to the lack of randomness in the 'S' boxes. Second was the publication of the paper New Directions in Cryptography by Whitfield Diffie and Martin Hellman. This paper introduced a radically new method of distributing cryptographic keys, which went far toward solving one of the fundamental problems of cryptography, key distribution. It has become known as Diffie-Hellman key exchange. The article also stimulated the almost immediate public development of a new class of enciphering algorithms, the asymmetric key algorithms. Prior to that time, all useful modern encryption algorithms had been symmetric key algorithms, in which the same cryptographic key is used with the underlying algorithm by both the sender and the recipient who must both keep it secret. All of the electromechanical machines used in WWII were of this logical class, as were the Caesar and Atbash cyphers and essentially all cypher and code systems throughout history. The 'key' for a code is, of course, the codebook, which must likewise be distributed and kept secret. Of necessity, the key in every such system had to be exchanged between the communicating parties in some secure way prior to any use of the system (the term usually used is 'via a secure channel') such as a trustworthy courier with a briefcase handcuffed to a wrist, or face-to-face contact, or a loyal carrier pigeon. This requirement rapidly becomes unmanageable when the number of participants increases beyond some (very!) small number, or when (really) secure channels aren't available for key exchange, or when, as is sensible crypto practice keys are changed frequently. In particular, a separate key is required for each communicating pair if no third party is to be able to decrypt their messages. A system of this kind is also known as a private key, secret key, or conventional key cryptosystem. D-H key exchange (and succeeding improvements) made operation of these systems much easier, and more secure, than had ever been possible before. In contrast, with asymmetric key encryption, there is a pair of mathematically related keys for the algorithm, one of which is used for encryption and the other for decryption. Some, but not all, of these algorithms have the additional property that one of the keys may be made public since the other cannot be (by any currently known method) deduced from the 'public' key. The other key in these systems is kept secret and is usually called, somewhat confusingly, the 'private' key. An algorithm of this kind is known as a public key / private key algorithm, although the term asymmetric key cryptography is preferred by those who wish to avoid the ambiguity of using that term for all such algorithms, and to stress that there are two distinct keys with different secrecy requirements. As a result, for those using such algorithms, only one key pair is now needed per recipient (regardless of the number of senders) as possession of a recipient's public key (by anyone whomsoever) does not compromise the 'security' of messages so long as the corresponding private key is not known to any attacker (effectively, this means not known to anyone except the recipient). This unanticipated, and quite surprising, property of some of these algorithms made possible, and made practical, widespread deployment of high quality crypto systems which could be used by anyone at all. Which in turn gave government crypto organizations worldwide a severe case of heartburn; for the first time ever, those outside that "fraternity" had access to cryptography that wasn't readily breakable by the 'snooper' side of those organizations. Considerable controversy, and conflict, began immediately. It has not yet subsided. In the US, for example, exporting "strong" cryptography remains illegal; cryptographic methods and techniques are classified as munitions. Until 2001 'strong' crypto was defined as anything using keys longer than 40 bits—the definition was relaxed thereafter. (See S Levy's "Crypto" for a journalistic account of the policy controversy in the US). Note, however, that it has NOT been proven impossible, for any of the good public/private asymmetric key algorithms, that a private key (regardless of length) can be deduced from a public key (or vice versa). Informed observers believe it to be currently impossible (and perhaps forever impossible) for the 'good' asymmetric algorithms; no workable 'companion key deduction' techniques have been publicly shown for any of them. Note also that some asymmetric key algorithms have been quite thoroughly broken, just as many symmetric key algorithms have. There is no special magic attached to using algorithms which require two keys. In fact, some of the well respected, and most widely used, public key / private key algorithms can be broken by one or another cryptanalytic attack and so, like other encryption algorithms, the protocols within which they are used must be chosen and implemented carefully to block such attacks. Indeed, "all" can be broken if the key length used is short enough to permit practical brute force key search; this is inherently true of all encryption algorithms using keys, including both symmetric and asymmetric algorithms. This is an example of the most fundamental problem for those who wish to keep their communications secure; they must choose a crypto system (algorithms + protocols + operation) that resists all attack from any attacker. There being no way to know who those attackers might be, nor what resources they might be able to deploy, nor what advances in cryptanalysis (or its associated mathematics) might in future occur, users may ONLY do the best they know how, and then hope. In practice, for well designed / implemented / used crypto systems, this is believed by informed observers to be enough, and possibly even enough for all(?) future attackers. Distinguishing between well designed / implemented / used crypto systems and crypto trash is another, quite difficult, problem for those who are not themselves expert cryptographers. It is even quite difficult for those who are. Revision of modern history. In recent years public disclosure of secret documents held by the UK government has shown that asymmetric key cryptography, D-H key exchange, and the best known of the public key / private key algorithms (i.e., what is usually called the RSA algorithm), all seem to have been developed at a UK intelligence agency before the public announcement by Diffie and Hellman in '76. GCHQ has released documents claiming that they had developed public key cryptography before the publication of Diffie and Hellman's paper. Various classified papers were written at GCHQ during the 1960s and 1970s which eventually led to schemes essentially identical to RSA encryption and to Diffie-Hellman key exchange in 1973 and 1974. Some of these have now been published, and the inventors (James Ellis, Clifford Cocks, and Malcolm Williamson) have made public (some of) their work. 



A transposition cipher encodes a message by reordering the plaintext in some definite way. Mathematically, it can be described as applying some sort of bijective function. The receiver decodes the message using the reordering in the opposite way, setting the ordering right again. Mathematically this means using the inverse function of the original encoding function. For example, to encrypt the sentence "A simple kind of transposition cipher writes the message into a rectangle by rows and reads it out by columns," we could use the following rectangle:  Asimplekin  doftranspo  sitionciph  erwritesth  emessagein  toarectang  lebyrowsan  dreadsitou  tbycolumns Then the encrypted text would be "Adsee tldts oirmo erbif tweab eymti rsrya cproi serdo lanta cosle ncegt wiuks iseas tmipp tinao nnohh ngnus." This cipher is often complicated by permuting the rows and columns, as in columnar transposition. Columnar transposition. The standard columnar transposition consists of writing the key out as column headers, then writing the message out in successive rows beneath these headers (filling in any spare spaces with nulls), finally, the message is read off in columns, in alphabetical order of the headers. For example suppose we have a key of 'ZEBRAS' and a message of 'WE ARE DISCOVERED. FLEE AT ONCE'. We start with: Then read it off as:  EVLNE ACDTK ESEAQ ROFOJ DEECU WIREE To decipher it, the recipient has to work out the column lengths by dividing the message length by the key length. Then he can write the message out in columns again, then re-order the columns by reforming the key word. Double transposition. A single columnar transposition could be attacked by guessing possible column lengths, writing the message out in its columns (but in the wrong order, as the key is not yet known), and then looking for possible anagrams. Thus to make it stronger, a double transposition was often used. This is simply a columnar transposition applied twice, with two different keys of different (preferably relatively prime) length. Double transposition was generally regarded as the most complicated cipher that an agent could operate reliably under difficult field conditions. It was in actual use at least as late as World War II (e.g. poem code). Grille. Another type of transpositional cipher uses a grille. This is a square piece of cardboard with holes in it such that each cell in the square appears in no more than one position when the grille is rotated to each of its four positions. Only grilles with an even number of character positions in the square can satisfy this requirement. As much message as will fit in the grille is written, then it is turned to another position and more message is written. Removing the cardboard reveals the cyphertext. The following diagram shows the message "JIM ATTACKS AT DAWN" encoded using a 4x4 grille. The top row shows the cardboard grille and the bottom row shows the paper underneath the grille at five stages of encoding: After the letters in the message have all been written out, the ciphertext can be read from the paper: "JKDT STAA AIWM NCAT". The sender and receiver must agree on the initial orientation of the grille, the direction to rotate the grille, the order in which to use the spaces on the grille, and the order in which to read the ciphertext characters from the paper. 

A Caesar cipher (also known as a shift cipher) is a substitution cipher in which the cipher alphabet is merely the plain alphabet rotated left or right by some number of positions. For instance, here is a Caesar cipher using a right rotation of three places:  Plain: ABCDEFGHIJKLMNOPQRSTUVWXYZ  Cipher: XYZABCDEFGHIJKLMNOPQRSTUVW To encipher a message, simply look up each letter of the message in the "plain" line and write down the corresponding letter in the "cipher" line. To decipher, do the reverse. Because this cipher is a group, multiple encryptions and decryptions provide NO additional security against any attack, including brute-force. The Caesar cipher is named for Julius Caesar, who allegedly used it to protect messages of military significance. It was secure at the time because Caesar's enemies could often not even read plaintext, let alone ciphertext. But since it can be very easily broken even by hand, it has not been adequate for secure communication for at least a thousand years since the Arabs discovered frequency analysis and so made all simple substitution cyphers almost trivially breakable. An ancient book on cryptography, now lost, is said to have discussed the use of such cyphers at considerable length. Our knowledge is due to side comments by other writers, such as Suetonius. Indeed, the Caesar cypher is much weaker than the (competently done) random substitution ciphers used in newspaper cryptogram puzzles. The most common places Caesar ciphers are found today are in children's toys such as secret decoder rings and in the ROT13 cipher on Usenet (which, of course, is meant to be trivial to decrypt)... 

In the field of cryptanalysis, frequency analysis is a methodology for "breaking" simple substitution ciphers, not just the Caesar cipher but all monoalphabetic substitution ciphers. These ciphers replace one letter of the plaintext with another to produce the cyphertext, and any particular letter in the plaintext will always, in the simplest and most easily breakable of these cyphers, turn into the same letter in the cypher. For instance, all E's will turn into X's. Frequency analysis is based on the fact that certain letters, and combinations of letters, appear with characteristic frequency in essentially all texts in a particular language. For instance, in the English language E is very common, while X is not. Likewise, ST, NG, TH, and QU are common combinations, while XT, NZ, and QJ are exceedingly uncommon, or "impossible". Given our example of all E's turning into X's, a cyphertext message containing lots of X's already seems to suggest one pair in the substitution mapping. In practice the use of frequency analysis consists of first counting the frequency of cypher text letters and then assigning "guessed" plaintext letters to them. Many letters will occur with roughly the same frequency, so a cypher with X's may indeed map X onto R, but could also map X onto G or M. But some letters in every language using letters will occur more frequently; if there are more X's in the cyphertext than anything else, it's a good guess for English plaintext that X stands for E. But T and A are also very common in English text, so X might be either of them. It's very unlikely to be a Z or Q which aren't common in English. Thus the cryptanalyst may need to try several combinations of mappings between cyphertext and plaintext letters. Once the common letters are 'solved', the technique typically moves on to pairs and other patterns. These often have the advantage of linking less commonly used letters in many cases, filling in the gaps in the candidate mapping table being built. For instance, Q and U nearly always travel together in that order in English, but Q is rare. Frequency analysis is extremely effective against the simpler substitution cyphers and will break astonishingly short ciphertexts with ease. This fact was the basis of Edgar Allan Poe's claim, in his famous newspaper cryptanalysis demonstrations in the middle 1800's, that no cypher devised by man could defeat him. Poe was overconfident in his proclamation, however, for polyalphabetic substitution cyphers (invented by Alberti around 1467) defy simple frequency analysis attacks. The electro-mechanical cypher machines of the first half of the 20th century (e.g., the Hebern? machine, the Enigma, the Japanese Purple machine, the SIGABA, the Typex, ...) were, if properly used, essentially immune to straightforward frequency analysis attack, being fundamentally polyalphabetic cyphers. They were broken using other attacks. Frequency analysis was first discovered in the Arab world, and is known to have been in use by about 1000 CE. It is thought that close textual study of the Koran first brought to light that Arabic has a characteristic letter frequency which can be used in cryptoanalysis. Its use spread, and was so widely used by European states by the Renaissance that several schemes were invented by cryptographers to defeat it. These included use of several alternatives to the most common letters in otherwise monoalphabetic substitution cyphers (i.e., for English, both X and Y cyphertext might mean plaintext E), use of several alphabets—chosen in assorted, more or less, devious ways (Leon Alberti seems to have been the first to propose this), culminating in such schemes as using only pairs or triplets of plaintext letters as the 'mapping index' to cyphertext letters (e.g., the Playfair cipher invented by Charles Wheatstone in the mid 1800s). The disadvantage of all these attempts to defeat frequency counting attacks is that it increases complication of both encyphering and decyphering, leading to mistakes. Famously, a British Foreign Secretary is said to have rejected the Playfair cipher because, even if school boys could learn it as Wheatstone and Playfair had shown, 'our attaches could never learn it!'. Frequency analysis requires a basic understanding of the language of the plaintext, as well as tenacity, some problem solving skills, and considerable tolerance for extensive letter bookkeeping. Neat handwriting also helps. During WWII, both the British and Americans recruited codebreakers by placing crossword puzzles in major newspapers and running contests for who could solve them the fastest. Several of the cyphers used by the Axis were breakable using frequency analysis (e.g., the 'consular' cyphers used by the Japanese). Mechanical methods of letter counting and statistical analysis (generally IBM card machinery) were first used in WWII. Today, the hard work of letter counting and analysis has been replaced by the tireless speed of the computer, which can carry out this analysis in seconds. No mere substitution cypher can be thought credibly safe in modern times. The frequency analysis method is neither necessary nor sufficient to solve ciphers. Historically, cryptanalysts solved substitution ciphers using a variety of other analysis methods long before and after the frequency analysis method became well known. Some people even question why the frequency analysis method was considered useful for such a long time. However, modern cyphers are not simple substitution cyphers in any guise. They are much more complex than WWII cyphers, and are immune to simple frequency analysis, and even to advanced statistical methods. The best of them must be attacked using fundamental mathematical methods not based on the peculiarities of the underlying plaintext language. See Cryptography/Differential cryptanalysis or Cryptography/Linear cryptanalysis as examples of such techniques. 

Before we can start any kind of examination of the abilities of the Bourne Shell and how you can tap into its power, we have to cover some basic ground first: we have to discuss how to enter commands into the shell for execution by that shell. The easy way: the interactive session. Taking another look at what you've probably already seen. If you have access to a Unix-based machine (or an emulator on another operating system), you've probably been using the Bourne Shell -- or one of its descendants -- already, possibly without realising. Surprise: you've been doing shell scripting for a while already! In your Unix environment, go to a terminal; either a textual logon terminal, or a terminal-in-a-window if you're using the X Window System (look for something called "xterm" or "rxvt" or just "terminal", if you have actually not ever done this yet). You'll probably end up looking at a screen looking something like this:  Ben_Tels:Local_Machine:~&gt;_ or  The admin says: everybody, STOP TRYING TO CRASH THE SYSTEM  Have a lot of fun!  bzt:Another_Machine:~&gt;_ or even something as simple as  $_ That's it. That's your shell: your direct access to everything the system has to offer. Using the shell in interactive mode. Specifically, the program you accessed a moment ago is your shell, running in "interactive mode": the shell is running in such a way that it displays a prompt and a cursor (the little, blinking line) and is waiting for you to enter a command for it to execute. You execute commands in interactive mode by typing them in, followed by a press of the Enter key. The shell then translates your command to something the operating system understands and passes off control to the operating system so that it can actually carry out the task you have sent it. You'll notice that your cursor will disappear momentarily while the command is being carried out, and you cannot type anymore (at this point, the Bourne Shell program is no longer in control of your terminal -- the other program that you started by executing your command is). At some point the operating system will be finished working on your command and the shell will bring up a new prompt and the cursor as well and will then start waiting again for you to enter another command. Give it a try: type the command After a short time, you'll see a list of files in the working directory (the directory that your shell considers the "current" directory), a new prompt and the cursor. This is the simplest way of executing shell commands: typing them in one at a time and waiting for each to complete in order. The shell is used in this way very often, both to execute commands that belong to the Bourne Shell programming language and simply to start running other programs (like the ls program from the example above). A useful tidbit. Before we move on, we'll mention two useful key combinations when using the shell: the command to interrupt running programs and shell commands and the command to quit the shell (although, why you would ever want to "stop" using the shell is beyond me...). To interrupt a running program or shell command, hit the Control and C keys at the same time. We'll get back to what this does exactly in a later chapter, but for now just remember this is the way to interrupt things. To quit the shell session, hit Control+d. This key combination produces the Unix end-of-file character -- we'll talk more later about why this also terminates your shell session. Some modern shells have disabled the use of Control+d in favor of the "exit" command (shame on them). If you're using such a shell, just type the word "exit" (like with any other command) and press Enter (from here on in, I'll leave the "Enter" out of examples). The only slightly less easy way: the script. As we saw in the last section, you can very easily execute shell commands for all purposes by starting an interactive shell session and typing your commands in at the prompt. However, sometimes you have a set of commands that you have to repeat regularly, even at different times and in different shell sessions. Of course, in the programming-centric environment of a Unix system, you can write a program to get the same result (in the C language for instance). But wouldn't it be a lot easier to have the convenience of the shell for this same task? Wouldn't it be more convenient to have a way to replay a set of commands? And to be able to compose that set as easily as you can write the single commands that you type into the shell's interactive sessions? The shell script. Fortunately, there "is" such a way: the Bourne Shell's "non-interactive" mode. In this mode, the shell doesn't have a prompt or wait for your commands. Instead, the shell reads commands from a text file (which tells the shell what to do, kind of like an actor gets commands from a script -- hence, shell script). This file contains a sequence of commands, just as you would enter them into the interactive session at the prompt. The file is read by the shell from top to bottom and commands are executed in that order. A shell script is very easy to write; you can use any text-editor you like (or even any wordprocessor or other editor, as long as you remember to save your script in plain text format). You write commands just as you would in the interactive shell. And you can run your script the moment you have saved it; no need to compile it or anything. Running a shell script. To run a shell script (to have the shell read it and execute all the commands in the script), you enter a command at an interactive shell prompt as you would when doing anything else (if you're using a graphical user interface, you can probably also execute your scripts with a click of the mouse). In this case, the program you want to start is the shell program itself. For instance, to run a script called "MyScript", you'd enter this command in the interactive shell (assuming the script is in your working directory): Starting the shell program from inside the shell program may sound weird at first, but it makes perfect sense if you think about it. After all, you're typing commands in an "interactive mode" shell session. To run a script, you want to start a shell in "non-interactive mode". That's what's happening in the above command. You'll note that the Bourne Shell executable takes a single parameter in the example above: the name of the script to execute. If you happen to be using a POSIX 1003.1-compliant shell, you can also execute a single command in this new, non-interactive session. You have to use the -c command-line switch to tell the shell you're passing in a command instead of the name of a script: We'll get to why you would want to do this (rather than simply enter your command directly into the interactive shell) a little further down. There is also another way to run a script from the interactive shell: you type the execute command (a single period) followed by the name of the script: The difference between that and using the "sh" command is that the "sh" command starts a new process and the execute command does not. We'll look into this (and its importance) in the next section. By the way, this notation with the period is commonly referred to as "sourcing" a script. Running a shell script the other way. There is also another way to execute a shell script, by making more direct use of a feature of the Unix operating system: the executable mode. In Unix, each and every file has three different permissions (read, write and execute) that can be set for three different entities: the user who owns the file, the group that the file belongs to and "the world" (everybody else). Give the command in the interactive shell to see the permissions for all files in the working directory (the column with up to nine letters, r, w and x for read write and execute, the first three for the user, the middle ones for the group, the right ones for the world). Whenever one of those entities has the "execute" permission, that entity can simply run the file as a program. To make your scripts executable by everybody, use the command as in You can then execute the script with a simple command like so (assuming it is in a directory that is in your PATH, the directories that the shell looks in for programs when you don't tell it exactly where to find the program): If this fails then the current directory is probably not in your PATH. You can force the execution of the script using At this command, the operating system examines the file, places it in memory and allows it to run like any other program. Of course, not every file makes "sense" as a program; a binary file is not necessarily a set of commands that the computer will recognize and a text file cannot be read by a computer at all. So to make our scripts run like this, we have to do something extra. As we mentioned before, the Unix operating system starts by examining the program. If the program is a text file rather than a binary one (and cannot simply be executed), the operating system expects the first line of the file to name the interpreter that the operating system should start to interpret the rest of the file. The line the Unix operating system expects to find looks like this: In our case, the following line should work pretty much everywhere: The Bourne Shell executable, to be found in the bin directory, which is right under the top of the filesystem tree. For example: Executing shell scripts like this has several advantages. First it's less cumbersome than the other notations (it requires less typing). Second, it's an extra safety if you're going to pass your scripts around to others. Instead of relying on them to have the right shell, you can simply specify which shell they should use. If Bourne Shell is enough, that's what you ask for. If you absolutely need "ksh" or "bash", you specify that instead (mind you, it's not foolproof — other people can ignore your interpreter specification by running your script with one of the other commands that we discussed above, even if the script probably won't work if they do that). Just as a sidenote, Unix doesn't limit this trick to shell scripts. Any script interpreter that expects its scripts to be plain-text can be specified in this way. You can use this same trick to make directly executable Perl scripts or Python, Ruby, etc. scripts as well as Bourne Shell scripts. Note also that with the distributions using bash as their default shell, you can use the #!/bin/sh shebang and have typical bash syntax in your script. It will work. But for the same script to work with a distribution not using bash as its default shell (as example Debian), you will have to modify the script or to change its shebang to #!/bin/bash. A little bit about Unix and multiprocessing. Why you want to know about multiprocessing. While this is not directly a book about Unix, there are some aspects of the Unix operating system that we must cover to fully understand why the Bourne Shell works the way it does from time to time. One of the most important aspects of the Unix operating system – in fact, the main aspect that sets it apart from all other main-stream operating systems – is that the Unix Operating System is and always has been a multi-user, multi-processing operating system (this in contrast with other operating systems like MacOS and Microsoft's DOS/Windows operating systems). The Unix OS was always meant to run machines that would be used simultaneously by several users, who would all want to run at least one but possibly several programs at the same time. The ability of an operating system to divide the time of a machine's processor among several programs so that it seems to the user that they are all running at the same time is called "multiprocessing". The Unix Operating System was designed from the core up with this possibility in mind and it has an effect on the way your shell sessions behave. Whenever you start a new process (by running a program, for instance) on your Unix machine, the operating system provides that process with its very own operating environment. That environment includes some memory for the process to play in and it can also include certain predefined settings for all processes. Whenever you run the shell program, it is running in its own environment. Whenever you start a new process from another process (for instance by issuing a command to your shell program in interactive mode), the new process becomes what is called a "child process" of the first process (the ls program runs as a child process of your shell, for instance). This is where it becomes important to know about multiprocessing and process interaction: a child process always starts with a "copy" of the environment of the parent process. This means two things: What does what. We have seen several ways of running a shell command or script. With respect to multiprocessing, they run in the following way: A useful thing to know: background processes. With the above, it may seem like multiprocessing is just a pain when doing shell scripting. But if that were so, we wouldn't "have" multiprocessing—Unix doesn't tend to keep things that aren't useful. Multiprocessing is a valuable tool in interacting with the rest of the system and one that you can use to work more efficiently. There are many books available on the benefits of multiprocessing in program development, but from the point of view of the Bourne Shell user and scripter the main one is the ability to hand off control of a process to the operating system "and still keep on working while that child process is running". The way to do this is to run your process as a "background process". Running a process as a background process means telling the operating system that you want to start a process, but that it should not attach itself to any of the interactive devices (keyboard, screen, etc.) that its parent process is using. And more than that, it also tells the operating system that the request to start this child process should return immediately and that the parent process should then be allowed to continue working without having to wait for its child process to end. This sounds complicated, but you have to keep in mind that this ability is completely ingrained in the Unix operating system and that the Bourne Shell was intended as an easy interface to the power of Unix. In other words: the Bourne Shell includes the ability to start a child process as a simple command of its own. Let's demonstrate how to do this and how useful the ability is at the same time, with an example. Give the following (rather pointless but still time consuming) command at the prompt: We'll get into what this says in later chapters; for now, it's enough to know that this command asks the system for the date and time and writes the result to a file named "scriptout". Since it then repeats this process 10000 times, it may take a little time to complete. Now give the following command: You'll notice that you can immediately resume using the shell (if you don't see this happening, hit Control+C and check that you have the extra ampersand at the end). After a while the background process will be finished and the scriptout file will contain another 10000 time reads. The way to start a background process in Bourne Shell is to append an ampersand (&amp;) to your command. Remarks. Actually, you can force a child process here as well -- we'll see how when we talk about command grouping 

A truing jig, aka a wheel truing stand, facilitates the precise truing or building of a wheel. It consists of fork, similar to a bicycle's front fork, and callipers which allow precise measurement of the position of the rim. It is possible to make your own truing jib. Start by getting an old bicycle fork and finding a way of supporting it. The easiest way is to find a sturdy box, and drilling a hole big enough to accommodate the fork column. The hole should be tight enough to firmly hold the fork and the wheel which will be supported. The fork is now upside down. It is possible to fashion some kind of true indicator or caliper using a nail or pencil attached on one leg of the fork. A metal strap that can be attached with a screw as a pencil or nail holder works best. I try to avoid using tape, string or rubber band since they have a tendency to move which could affect an accurate reading of the wheel's straightness. Straight wheels tend to turn smoothly without deviation when viewed from the front. I start by bringing the nail as closely as possible to the wheel rim. Turning the wheel by hand, warped wheels will occasionally touch the rim or show a large gap. Wheel spokes work by providing enough pull or release on both sides of the rim so that the rim remains straight. By loosening one spoke, the tension is strengthened on the other side and vice-versa. Ideally the spokes on both sides should be even to maintain optimum rim straightness. A sequence of tightening and loosening both sides of the spokes where necessary is the key to rim straightness. 

A brute force attack against a cipher consists of breaking a cipher by trying all possible keys. Statistically, if the keys were originally chosen randomly, the plaintext will become available after about half of the possible keys are tried. As we discuss in ../Basic Design Principles/, the underlying assumption is, of course, that the cipher is known. Since A. Kerckhoffs first published it, a fundamental maxim of cryptography has been that security must reside only in the key. As Claude E. Shannon said a few decades later, 'the enemy knows the system'. In practice, it has been excellent advice. As of the year 2002, symmetric ciphers with keys 64 bits or fewer are vulnerable to brute force attacks. DES, a well respected symmetric algorithm which uses 56-bit keys, was broken by an EFF project in the late 1990s. They even wrote a book about their exploit—"Cracking DES", O'Reilly and Assoc. The EFF is a non-profit cyberspace civil rights group; many people feel that well-funded organisations like the NSA can successfully attack a symmetric key cipher with a 64-bit key using brute force. This is surely true, as it has been done publicly. Many observers suggest a minimum key length for symmetric key algorithms of 128 bits, and even then it is important to select a secure algorithm. For instance, many algorithms can be reduced in effective keylength until it is computationally feasible to launch a brute force attack. AES is recommended for use until at least 2030. The situation with regard to asymmetric algorithms is much more complicated and depends on the individual algorithm. Thus the currently breakable key length for the RSA algorithm is at least 768 bits (broken publicly since 2009), but for most elliptic curve asymmetric algorithms, the largest currently breakable key length is believed to be rather shorter, perhaps as little as 128 bits or so. A message encrypted with a 109 bit key by an elliptic curve encryption algorithm was publicly broken by brute force key search in early 2003. As of 2015, a minimum key length of 224 bits is recommended for elliptic curve algorithms, and 2048 bits for such other asymmetric key algorithms as RSA (asymmetric key algorithms that rely on complex mathematical problems for their security always will need much larger keyspaces as there are short-cuts to cracking them, as opposed to direct brute-force). Common Brute Force Attacks. The term "brute force attacks" is really an umbrella term for all attacks that exhaustively search through all possible (or likely) combinations, or any derivative thereof. Dictionary Attack. A dictionary attack is a common password cracking technique, relying largely on the weak passwords selected by average computer users. For instance, if an attacker had somehow accessed the hashed password files through various malicious database manipulations and educated searching on an online store, he would then write a program to hash one at a time all words in a dictionary (of, for example any or all languages and common derivative passwords), and compare these hashes to the real password hashes he had obtained. If the hashes match, he has obtained a password. Pre-Computation Dictionary Attack. The simple dictionary attack method quickly becomes far too time-consuming with any large number of password hashes, such as an online database would yield. Thus, attackers developed the method of pre-computation. In this attack, the attacker has already hashed his entire suite of dictionaries, and all he need do is compare the hashes. Additionally, his task is made easier by the fact that many users will select the same passwords. To prevent this attack, a database administrator must attach unique 32-bit salts to the users passwords before hashing, thus rendering precompution useless. The ../Breaking Hash Algorithms/ chapter of this books goes into more detail on attacks that specifically apply to hashed password files. Responses to Brute Force Attacks. There are a number of ways to mitigate brute force attacks. For example: The ../Secure Passwords/ chapter of this book goes into more detail on mitigations and other responses that specifically apply to hashed password files. 

The Data Encryption Standard (DES) was a widely-used algorithm for encrypting data. It was developed by IBM under the name Lucifer, and was submitted to NBS in response to a 1973 solicitation for better cryptosystems. The US National Institute of Standards and Technology with help from the National Security Agency took IBM's design and made some changes; DES was adopted as a standard in January 1977. DES is a product block encryption algorithm (a cipher) in which 16 iterations, or rounds, of the substitution and transposition (permutation) process are cascaded. The block size is 64 bits, so that a 64-bit block of data (plaintext) can be encrypted into a 64-bit ciphertext. The key, which controls the transformation, also consists of 64 bits. Only 56 of these, however, are at the user's disposal; the remaining eight bits are employed for checking parity. The actual key length is therefore 56 bits. Subsets of the key bits are designated K1, K2, etc., with the subscript indicating the number of the round. The cipher function (substitution and transposition) that is used with the key bits in each round is labeled f. At each intermediate stage of the transformation process, the cipher output from the preceding stage is partitioned into the 32 leftmost bits, Li, and the 32 rightmost bits, Ri. Ri is transposed to become the left-hand part of the next higher intermediate cipher, Li+1. The right-hand half of the next cipher, Ri+1, however, is a complex function of the key and of the entire preceding intermediate cipher. The essential feature to the security of the DES is that f involves a very special nonlinear substitution—i.e., f(A) + f(B) does not equal f(A + B)--specified by the Bureau of Standards? in tabulated functions known as S-boxes. This operation results in a 32-bit number, which is logically added to Ri to produce the left-hand half of the new intermediate cipher. This process is repeated, 16 times in all. To decrypt a cipher, the process is carried out in reverse order, with the 16th round being first. The DES algorithm lends itself to integrated-chip implementation. By 1984 the Bureau of Standards had certified over 35 LSI- and VLSI-chip implementations of the DES, most on single 40-pin chips, some of which operate at speeds of several million bits per second. When the cipher was first released, the design criteria for the S-boxes was not released. With the National Security Agency's involvement in the design of the S-boxes, most security researchers were wary of DES, and there was the widespread fear that the modifications of the NSA were intended to weaken the cipher. In 1990 with the independent discovery and open publication by Biham and Shamir of differential cryptanalysis, it turned out that at least some of the wariness was uncalled for. After the publication of this paper, the IBM personnel involved in the designed publically stated that the main factor in the design was to strengthen them against differential cryptanalysis. The secrecy behind the design criteria at the time appears to have been due to the fact that the technique was not known to the public at the time. Notably, DES is theoretically vulnerable to a technique discovered later by Matsui, linear cryptanalysis. It is unknown whether the NSA was aware of linear cryptanalysis at the time DES was finalized, but most knowledgeable observers think not. Don Coppersmith, one of DES's designers at IBM, has stated that IBM itself was not aware of linear cryptanalysis at that time. Because the key length is only 56 bits, DES can be, and has been, broken by the brute force attack method of running through all possible keys. It is believed that one of the reasons this reduced key length was chosen was that NSA in the mid-'70s possessed enough computer power to brute force break keys of this length. In the years since, computer hardware progress has been such that most anyone now can have sufficient computational capacity. The EFF, a cyberspace civil rights group (with neither much funding nor personnel), did it in a little more than 2 days' search at about the same time at least one attorney from the US Justice Department was publicly announcing that DES was and would remain unbreakable. The most obvious way of improving the security of DES is to encrypt the data multiple times with different keys. Double encrypting data with DES does not add much security as it is vulnerable to meet in the middle attacks. Going one step about this, many former DES users now use Triple DES (3DES) which was described and analyzed by one of DES's patentees (see FIPS 46-3); it involves DES encryption of each data block three times with different keys. 3DES is widely regarded as adequately secure for now, though it is quite slow. Note, however, that there are several ways to use DES three times; only one of those is Tuchman's 3DES. After another, long delayed competition, (NIST) has selected a new cipher, the Advanced Encryption Standard (AES) to replace DES (fall -'01). AES was submitted by its designers under the name Rijndael. Implementations: http://www.codeproject.com/KB/cs/NET_Encrypt_Decrypt.aspx (C#, Xinwen Cheng) &lt;br&gt; http://frank.anemaet.nl/crypto/DES/ (Java implementation, Frank Anemaet) &lt;br&gt; http://www.tero.co.uk/des/ (Javascript implementation, Paul Tero) 

=Contributors= The Chinese Wikibook was started 2003 December 13. Below is a list of users who have contributed greatly to the authoring of this Wikibook. Please add your username if you have made substantial additions and/or revisions to this textbook. Use * to add a name. , and all made substantial contributions to the on , which are used in our lessons. , also of Wikicommons, contributed the first audio samples used in this Wikibook. In addition, the authors would like to thank the development team in relation with the Wikimedia Foundation and its affiliates, without whom our text could not be so accessible. 

Goals. Book. The goal of this book is to provide a comprehensive coverage of eXtensible Markup Language (XML) in a textbook format. This book is written and edited by students for students. Each student who uses the book should improve its quality by correcting errors, adding exercises, adding examples, starting new chapters and so forth. Chapters 2 through 6 take the perspective that an XML schema is a representation of a data model, and thus these chapters deal with mapping the complete range of relationships that occur between entities. As you learn how to convert each type of relationship into a schema, other aspects of XML are introduced. For example, stylesheets are initially introduced in Chapter 2 and progressively more stylesheet features are added in Chapters 3 through 6. Consolidation chapters (e.g., Chapter 7 "Data schemas") bring together the material covered across previous chapters; in this case, Chapters 2 through 6. This means students see key skills twice: once in the context of gradually developing their broad understanding of XML and then again in the specific context of one dimension of XML. Application chapters cover particular uses of XML (e.g., SVG for scalable vector graphics) to give the reader examples of the use of XML to solve particular types of problems. This part of the book is expected to grow as the use of XML extends. Project. Professors typically throw away their students’ projects at the end of the term. This is a massive waste of intellectual resources that can be harnessed for the betterment of many by creating an appropriate infrastructure. In our case, we use wiki technology as the infrastructure to create a free open content textbook. University students are an immense untapped global resource. They can be engaged in creating open textbooks if the right infrastructure is in place to sustain renewable student projects. This book is an example of how waste can be avoided. Software. To complete the exercises in the book and view the slides, you will need access to the following software (or a suitable alternative): 

 __NOEDITSECTION__ About This Wiki. This work is disjointly licensed under the GNU FDL, the CC BY SA, and Io's license. The authors of this tutorial are: And many others who wish to remain anonymous. Hello world. When learning to program, the first thing usually learned is how to write a program that prints "Hello, world!" to the console. It gives programmers a feel of how simple programs are structured. Io has many ways to write things to the console, so at times it can be confusing about which one to use. The one that's most often used is the write function. Enter the Io's interactive interpreter (by typing codice_1 at the console) and type the following:  write("Hello, world!\n") With the interactive interpreter, this should look like:  user@computer:~$ io  Io&gt; write("Hello, world!\n")  Hello, world!  ==&gt; nil Another way would be:  "Hello, world!\n" print Some of this makes sense. It looks like codice_2 is telling Io to write something to the console, specifically the stuff after it. The stuff after it is in quotes so it's some text, and it's what's going to be printed. But why are those parentheses there, what's with that codice_3 at the end of it, and why was there an arrow with the quoted text after it? The parentheses are telling codice_2 to use the thing within the parentheses, which is called an "argument". codice_2 is called a "function" (remember functions in math?). There might be multiple arguments within the parentheses separated by commas in a single function. The first and only argument for codice_2, namely codice_7, is a "string" (because it's a string of letters). The codice_3 at the end of the string is an "escape sequence". You can tell that because it begins with a backslash. This particular escape sequence signals a newline, the equivalent of pressing enter. If you just pressed enter in the middle of the string (or at least this type of string -- we'll get to that later), it would confuse Io. The thing where there's an arrow and then the string that was the argument of the codice_2 function is the value that codice_2 "return"ed. Every function returns a value. With some functions, it makes sense to return something (such as a function to add two numbers), but with other functions, it just returns something simple done with the arguments. In this case, codice_2 returns the special value codice_12, indicating it has nothing to return. Concatenate Strings. Use the double dot operator to concatenate strings. Or, you can use string interpolation, as shown below. Let's say codice_13 is a codice_14 instance, and it has both a first name and a last name:  Io&gt; olle  ==&gt; Object_0080B818 do(  appendProto(Object_00806B58)  lastname := "Jonsson"  title := "Developer"  firstname := "Olle" We want a method to show first and last name with a space in-between:  Io&gt; olle fullname := method(firstname .. " " .. lastname)  ==&gt; method(firstname ..(" ") ..(lastname)) or, using string interpolation:  Io&gt; olle fullname := method("#{firstname} #{lastname}" interpolate)  ==&gt; method("#{firstname} #{lastname}" interpolate) Now, calling codice_15 will yield codice_16. Looking at the returned method object, we can see how the Io interpreter creates regular message calls (the parentheses) for the method codice_17. So, this is the new codice_13 instance:  Io&gt; olle  ==&gt; Object_0080B818 do(  appendProto(Object_00806B58)  fullname := Block_007B5030  lastname := "Jonsson"  title := "Developer"  firstname := "Olle" Simple arithmetic. In Io, you can use arithmetic expressions and they will work correctly. Arithmetic in Io is just like functions, except it uses objects to help. Objects make it so that instead of typing something like +(1, 2), you can type 1+2. More about objects later. + is still a function, though, so it returns a value. This is very useful; it makes it so you can use Io as a simple calculator. Notes about what it's doing are put after two slashes (like this: //).  Io&gt; 1+2 //addition  ==&gt; 3  Io&gt; 4-5 //subtraction  ==&gt; -1  Io&gt; 7*3 //multiplication  ==&gt; 21  Io&gt; 3/6 //division  ==&gt; 0.50000  Io&gt; 2**3 //exponents  ==&gt; 8  Io&gt; 7%2 //remainder of 7/2 (technically, 7 mod 2)  ==&gt; 1 This follows normal order of operations (called "precedence") and parentheses can be used. As you'd expect, you can use codice_2 on numbers, but the newline isn't included. You have to use multiple arguments with codice_2 to print multiple things and then it returns all of those things put together. So to print codice_27, you would write:  write((1/3)**2, "\n") As you would expect, this prints 0.11111. It returns the string "0.11111\n" because with the codice_2 function, multiple arguments are converted to strings and joined together. Then that value is printed and returned. Variables. A "variable" is basically a word that stands for a value. They are somewhat like variables in mathematics, except in mathematics, variables could only be numbers and they were one letter long. In programming, functions, objects, strings, and numbers are all types of variables, but we haven't defined any yet. Many variables, such as write, 3, and "Hello, world!" can already be used, but only some of them can be changed. You can make your own variables using = and :=. Below are some examples of making and using variables:  x := 3  line := "\n"  write(x, line) Can you tell what's happening? The variable x is being set to 3 and the variable line is being set to "\n", which is equivalent to a newline. Then, the contents x and line are being written to the console. Since x is 3 and line is a newline, this prints 3 and then goes to the next line on the console. The function returns "3\n". In Io, there is a difference between creating and changing the value of variables. If a variable doesn't exist yet, you have to use :=, but if you've already given it an initial value using :=, you can use = for subsequent definitions. Here's an example:  x := 1  incrementor := 2  write("x is ", x, "\n")  x = x + incrementor  write("but when we added ", incrementor, " to it, it became ", x, "\n") It may be confusing that sometimes we use := and other times we use =, but you'll get used to it. If you want to, you can always use :=, but when we get into objects later, it will become very inconvenient to keep using :=, so you should probably start using = whenever appropriate. Programs. Until now, you've been simply been going to Io's interactive interpreter. This prevents you from making larger applications or writing things for others. If you want to write code to be reused, simply put it in a text file and run it with codice_29, where &lt;filename&gt; is the name of the file you use. Here's an example of using a file on Linux to store a program:  user@computer:~$ cat &gt; incrementor.io  x := 1  incrementor := 2  writeln("x is ", x)  x = x + incrementor  writeln("but when we added ", incrementor, " to it, it became ", x)  user@computer:~$ io incrementor.io  x is 1  but when we added 2 to it, it became 3  user@computer:~$ That was the last example we just did. If you noticed, I used a .io file extension. This is in no way mandated, it is merely a convention. Something happened differently this time then when we did it from a file, if you noticed. Unlike before, we didn't see what each function returned. Instead, we only saw what was explicitly output by codice_2. If you're on Linux or a similar system (such as Unix, Mac OS X, or Cygwin), you can make it so that your file can be run simply by typing codice_31 (where yourProgram.io is the name of your program), without needing to precede it with codice_1. This can be accomplished by putting a line of code at the beginning of your program and then giving it executable permissions. On Linux, that code is:  #! /usr/bin/env io It may differ for other systems. To give it executable permissions, simply type the following:  chmod +x yourProgram.io Again, this may be different on different systems. Once you do this, nothing will change about the actual execution of the code, but you can, for example, just double click on a source code file in a GUI and it will run. Writing functions. As I said before, functions are just another type of variable, so you can create them using := and change them using =. Functions themselves are created with a function called codice_33. Here's an example of a function:  add := method(a, b, //function to add 2 numbers  a + b  writeln(add(1, 2)) //writes 3  x := 1  writeln(add(3, x)) //writes 4  x = add(4, 5) //9  writeln(x) //writes 9  x = add(x, 1) //in effect increment x by 1  writeln(x) //writes 10 The function add takes two arguments, a and b, and then returns a + b. It would do exactly the same thing if, instead of a and b, we used the variable names 'this' and 'that', everywhere that a and b were used. In that case, the function would be written as:  add := method(this, that,  this + that What a function really does is take a list of arguments; then it assigns the arguments to "local variable"s, which are specified when you make the function. Local variables stay in the function you're working in and you can't get to them outside of the function. In the most recent case, those local variables are 'this' and 'that', so the variable 'this' will take the first argument, and the variable 'that' will take the second argument. Then it executes the rest of the function and returns the result. The function can have multiple lines of code and the last line will be the value returned. Here's an example:  examineArgs := method(this, that,  writeln("This is ", this, ".")  writeln("That is ", that, ".")  // usage:  examineArgs(3, 5)  /* writes:  This is 3.  That is 5. */  x := examineArgs("hi", "bye")  /* writes:  This is hi.  That is bye. */  write(x)  /* writes:  That is bye. */ Why was x set to "That is bye.\n"? Because that's what examineArgs returned. It returned the value of the last line, the value of write("That is ", that, ".\n"). That was set to "bye", so it printed the string "That is bye.\n" and returned the same string. If you want to return something before the function ends, you can use the return function. Unlike most functions, you don't need to put parentheses around the argument to call the return function. Here is an example of its usage:  returning := method(this, that,  writeln(this)  return this  writeln(that)  x := returning(1, 2)  //writes "1\n"  write(x) //writes 1 with no newline As you can see, the code after return isn't used at all. Right now, this doesn't seem very useful, but later, when we get into flow control, it will be used very often. Conditionals. Conditionals in Io are made using the if function:  if(a == 1) then(  writeln("a is one")  ) else(  writeln("a is not one") However, the preferred way to write this (without the need for then() and else() messages, so it is faster): 



Breaking the Caesar cipher is trivial as it is vulnerable to most forms of attack. The system is so easily broken that it is often faster to perform a brute force attack to discover if this cipher is in use or not. An easy way for humans to decipher it is to examine the letter frequencies of the cipher text and see where they match those found in the underlying language. Frequency analysis. By graphing the frequencies of letters in the ciphertext and those in the original language of the plaintext, a human can spot the value of the key but looking at the displacement of particular features of the graph. For example in the English language the frequencies of the letters Q,R,S,T have a particularly distinctive pattern. Computers can also do this trivially by means of an auto-correlation function. Brute force. As the system only has 25 non-trivial keys it is easy even for a human to cycle through all the possible keys until they find one which allows the ciphertext to be converted into plaintext. Known plaintext attack. If you have a message in both ciphertext and in plaintext it is trivial to find the key by calculating the difference between them. 

Computer languages come in many flavors. Generally there are high-level and low-level languages. The term "language level" is an attempt to convey information about how close the programming language is to the machine language of the particular hardware. Low-level languages are closer to the machine language, and high-level languages are further away from machine language, and close to natural languages (e.g., English). An extreme example of a low-level language is the machine language which consists of ones and zeros. It is extremely rare (difficult, and error prone), however, for people to write programs in machine language. The closest programming language to machine language is assembly language, which normally has one instruction, or statement, for each machine language instruction. An example of the kind of instruction that exists in assembly language is to copy the contents of a particular memory address into a machine register. One example of a high-level computer language might be , where the language is programmed by the user with specific words which can be found in the English dictionary. Any language (computer or natural) still must have certain structures, though. English sentences are made of English words which are combined to respect the grammar (i.e., the structure) of English. We understand a sentence if the grammar is (more or less) correct. However, just because a high-level language uses some words which seem to be natural English does not mean that the computer can directly understand English sentences as we would speak them to a person. A program must be written to respect the structure (grammar) of the computer language used, if the computer is to understand. Typically the grammar of a computer language is more rigid than the grammar of a natural language. 



This book deals with Engineering Thermodynamics, where concepts of thermodynamics are used to solve engineering problems. Engineers use thermodynamics to calculate the fuel efficiency of engines, and to find ways to make more efficient systems, be they rockets, refineries, or nuclear reactors. One aspect of "engineering" in the title is that a lot of the data used is empirical ("e.g." steam tables), since you won't find clean algebraic equations of state for many common working substances. Thermodynamics is the science that deals with transfer of heat and work. Engineering thermodynamics develops the theory and techniques required to use empirical thermodynamic data effectively. This course forms the foundation for the Heat Transfer course, where the "rate" and mechanisms of transmission of energy in the form of heat is studied. The concepts will be used in further courses in heat, Internal Combustion Engines, Refrigeration and Air Conditioning, and Turbomachines to name a few. Rigorous treatment of the molecular basis will be omitted, in favor of formulations most useful for developing intuition and understanding common technologies. Students of physics will want to pair this text with one on Statistical Mechanics. 

How to Find a Book depends on how much information you have about the book and whether you want to buy it or simply read it. Acquiring information about the book. Normally, title and name of the author of a book is enough to locate more information about it. You can then use online bookstores or online library catalogs to identify publisher, publication date, edition and ISBN; these items are sometimes required by libraries or bookstores to order the book. Occasionally, the author's name is misspelled, especially with foreign authors. In this case try searching for the title alone, until you find the proper spelling of the name. If you only have the book's ISBN, you should be aware of the fact that different editions as well as hardcover and soft-cover versions of a book have different ISBNs. It is possible that the version of the book with your ISBN is already out of print, while another version is still available. It is therefore best to use an online service to get the author and title information, and then proceed with that information. Locating a book is much more difficult if you want to educate yourself about a topic but don't have a particular author or title in mind. Library shelves are organized by subject; learn the system used by your library, go to the relevant shelve and browse the books about your topic. Some of this can also be done online: library catalogs allow subject and keyword searches and the Library of Congress permits browsing by call number, which can often locate closely related materials. For an effective subject search, you need to know which subject words are being used; this is contained in the "Library of Congress Subject Headings", a red five-volume work that many libraries have but can not be found online (most library subjects can be searched and browsed on ISBNdb.com's subject search). Most topics have one or two "standard references", the generally acknowledged best treatments of the subject. To find these, you need to talk to a specialist in the field. A Usenet posting to a relevant newsgroup, asking about good books, also often yields excellent results. Closely reading the user-supplied book reviews on amazon.com will also often point you to standard references: they are often only identified by author and everyone seems to know them. Acquiring the book. You have to decide whether you want to buy the book or just read it. If you want to buy it, you have to decide whether to buy it new or used (the only option if no version of it is in print, unless you can find a new after-market copy or a remaindered book). To buy used books, you can try a local used book store. These can also often help you buy books that they don't have. You may also go to a professional book finder with connections in the used-book markets. These used-book markets are now open to everyone with a web browser and a credit or debit card, and books can also be found via auction websites. The three best sites for used-books, which between them cover countless thousands of sellers are: Abebooks, eBay and Amazon (on an Amazon book entry look for the link to "new and used books", which takes one to many sellers of used books). A book in print can be ordered at any bookstore and can also be ordered online; several price comparison services exist. By ordering online, you typically save sales tax but you have to pay shipping and handling costs. Depending on what you need, you may be able to find a book at a public library or a good local college or university library. These can be used by non-students for free or for a small fee. If your reading interest is general, or if you want to read a popular novel, a public library is the place to go. Membership is usually free. If you cannot find the book on your own, you should visit an information desk because many libraries offer the ability to place "holds" on materials so you can get them when you return. At college and university libraries, Research librarians can help you locate books and journal articles of academic interest quickly. On campus, you can often use expensive databases such as LexisNexis for free. Libraries are connected by an Inter Library Loan network: if the book exists, and is held by a library willing to loan it, you will have it in as little as a week or so. Most libraries also offer an Inter-library loan service. This means that items that do not appear in their catalog can be requested from libraries elsewhere. Depending on the library and how far the book has to travel, there may be charges for this service. Libraries use WorldCat at OCLC in order to find books at other libraries. It is by far the most extensive database of library holdings. Some libraries allow access by their patrons; ask your librarian how to access WorldCat. Unsuccessful ABEbooks searches are directed to WorldCat as unsuccessful WorldCat searches are to ABEbooks. WorldCat attempts to restrict use of its data but such terms are not legally valid under US copyright law, since only the specific and creative original expression, if any, can be protected by copyright and the right to use public domain works cannot be restricted by contract (see Feist v. Rural and Assessment Technologies v. WIREdata) External links. Digital libraries. There are also digital libraries, where you can find books in electronic form (mostly scanned paper books). The most famous digital library is Project Gutenberg. 

Chemistry is Everywhere. The modern human experience places a large emphasis upon the material world. From the day of our birth to the day we die, we are frequently preoccupied with the world around us. Whether struggling to feed ourselves, occupying ourselves with modern inventions, interacting with other people or animals, or simply meditating on the air we breathe, our attention is focused on different aspects of the material world. In fact only a handful of disciplines—certain subsets of religion, philosophy, and abstract math—can be considered completely unrelated to the material world. Everything else is somehow related to chemistry, the scientific discipline which studies the properties, composition, and transformation of matter. Branches of Chemistry. Chemistry itself has a number of branches: Chemistry as a discipline is based on a number of other fields. Because it is a measurement-based science, math plays an important role in its study and usage. A proficiency in high-school level algebra should be all that is needed in this text, and can be obtained from a number of sources. Chemistry itself is determined by the rules and principles of physics. Basic principles from physics may be introduced in this text when necessary. Why Study Chemistry? There are many reasons to study chemistry. It is one pillar of the natural sciences necessary for detailed studies in the physical sciences or engineering. The principles of biology and psychology are rooted in the biochemistry of the animal world, in ways that are only now beginning to be understood. Modern medicine is firmly rooted in the chemical nature of the human body. Even students without long-term aspirations in science find beauty in the infinite possibilities that originate from the small set of rules found in chemistry. Chemistry has the power to explain everything in this world, from the ordinary to the bizarre. Why does iron rust? What makes propane such an efficient, clean burning fuel? How can soot and diamond be so different in appearance, yet so similar chemically? Chemistry has the answer to these questions, and so many more. Understanding chemistry is the key to understanding the world as we know it. This Book: General Chemistry. An introduction to the chemical world is set forth in this text. The units of study are organized as follows. « Begin Your Study of General Chemistry! » 

The MIPS microprocessor paradigm was created in 1981 from work done by J. L. Hennessy at Stanford University. Since that time, the MIPS paradigm has been so influential that nearly every modern-day processor family makes some use of the concepts derived from that original research. This book will discuss the MIPS architecture and (perhaps more importantly) MIPS assembly programming. Resources and Licensing.  __NOEDITSECTION__ 

alternate subtracting and adding numbers starting with 1 and increasing by 1 each time. 3,2,4,1,5,0,6,-1,7... 

&lt; Problem 3 2 4 1 4 ? Solution Alternate subtracting and multiplying, starting with 1 and increasing by 1 each time. 3,2,4,1,4,-1,-6... 

Introduction. Chess is an ancient strategy game that originated in India. It is played by two individuals on an 8×8 grid. The objective is to maneuver one's pieces so as to trap the opposing king in "checkmate". This book will cover the basic pieces of chess, before going on to some more advanced topics. The history of chess began in India during the Gupta Empire, where its early form in the 6th century CE was known as "chaturanga", which translates as "four divisions of the military" – infantry, cavalry, elephants, and chariotry, represented by the pieces that would evolve into the modern pawn, knight, bishop, and rook, respectively. In Sassanid Persia around 600 CE, the name became "shatranj", and the rules were developed further. Shatranj was taken up by the Muslim world after the Islamic conquest of Persia, with the pieces largely retaining their Persian names. In Spanish, "shatranj" was rendered as "ajedrez", in Portuguese as "xadrez", and in Greek as "zatrikion", but in the rest of Europe it was replaced by versions of the Persian "shāh" ("king"). 

General Chemistry Physical changes in chemistry include phase changes and anything else that changes the way that matter is arranged in space. Examples of physical changes are: Phase changes. The three basic phases of matter are solid, liquid, and gas. (Under certain unusual conditions matter can transform into a phase called plasma) A phase change is a change from one phase to another. The most common example is liquid water freezing into ice or evaporating into a gas. Phase changes result in different properties for the substance changing phases but the chemical identity of the material does not change. In the case of water, its molecules are always composed of two hydrogen atoms bonded to one oxygen atom. However, in frozen water the molecules are frozen into place in relation to each other in a structure called a crystal; in liquid water the molecules flow relative to each other; and water molecules in the gas phase are flying freely in space and seldom even contact each other. Phase changes are controlled by temperature and pressure. By lowering the temperature of a gas it can be condensed into a liquid and then even into a solid. Heating a solid melts it into a liquid and then, with further heat, into a gas. 

Appendix 3 - Voornaamwoorden ~ "pronouns". Like English, Dutch has pronouns. These can mark number, case, gender,politeness and emphasis. Pronouns can function either as substantives (nouns) or as adjectives. There is also a number of related adverbs that will be treated here. Adverbs are typically not considered pronouns in grammatical analysis, but they deserve mention when discussing the Dutch language because pronouns are often "replaced" by pronominal adverbs. Persoonlijke voornaamwoorden ~ "Personal pronouns". In this table personal pronouns are given in nominative, accusative and dative case. These cases signify the role the pronouns have in the sentence. For example: In "I am hitting you", "I" is nominative (subject) and "you" is accusative (object). Also words with a preposition are in accusative case ("you" in "I am looking at you"). Dative case is special and tells us something is indirect object, as "me" in "He gave me that" or "He built me a snowman" or, with a preposition, "He gave it to me". Remarks: Bezittelijke voornaamwoorden ~ "Possessive pronouns". Pssessive pronouns are essentially the adjectival forms of the personal pronouns. Remarks: Personal Adverb - "er". Dutch has a somewhat curious personal locative adverb er that replaces "het" and "ze" particularly in inanimate cases (i.e. for things more so than for persons). It occurs as the locative part of many pronominal adverbs, such as :erin, erdoor, ervan etc. but it can also be used independently: Notice that "er" is not considered the subject of these sentences ("koffie" and "mensen" are the subject resp.) Aanwijzende voornaamwoorden -- Demonstrative pronouns. Demonstrative pronouns are typically used as adjectives: they can also be used independently: They are more and more used to replace inanimate personal pronouns. Aanwijzende bijwoorden - demonstrative adverbs. Dutch has three demonstrative adverbs of time: One modal demonstrative adverb is common: Occasionally a more proximate one "zus" is used for contrast Two locative adverbs are in common use: Both of them are used as the locative part of demonstrative pronominal adverbs like: hierdoor, daarvan etc. A third adverb is less common: Betrekkelijke voornaamwoorden -- Relative pronouns. Zelfstandig- substantive. Without "antecedent": With inclusion of antecedent. There are a number of archaic forms that can be used with prepositions: As in English the genitives wiens and wier (whose) can be used in relative clauses referring to persons: In inanimate cases the relative pronominal adverb waarvan is virtually mandatory. Bijwoordelijk - adverbial. waar - where "Waar" is also used to form the relative pronominal adverbs like waarvan, waarvoor etc. that frequently replace relative pronouns. 

Powers of ten: 

The list below is partial list of colours/colors described in Vietnamese: Examples: 

Gesprek 3-1. Mam teaches her toddler, Jeroen to count: Leren 3 ~ Tellen van 1 tot 12. In Dutch, as in English, there are both ordinal and cardinal numbers, and number formation is similar in that the first twelve numbers are unique. Above twelve, numbers are formed by combination. For example, 15 is "vijftien" and 16 is "zestien". Other numbers will be the subject of more advanced lessons. Note in the table how ordinals are formed from the cardinals in Dutch by adding -de. 'Ten' becomes 'tenth' in English; tien become tiende in Dutch. As in English, there are several variants: "eerste", "derde", and "achtste". Remark: een is used both as an indefinite article ("a" or "an") and a number ("one"). One often puts accents on the e's when "one" is meant in case of ambiguity: één. In Dutch spelling such accents are allowed but only if otherwise ambiguity would arise. There is also a difference in pronunciation: /ən/ (-n) for the article and /e:n/ (ayn) for the number. Eerst en laatst. The ordinals are a special kind of adjectives. They always have the inflection -e. So, words like *zesd do not exist. The only exception is "eerst". As in English, it can be used as an adverb: Its opposite (antonym) is "laatst" as adverb and "laatste" as adjective: Grammatica 3-1 ~ Telling time (hours). Knowing the numbers from 1 to 12, you can now begin asking and telling time in Dutch. Gesprek 3-2. Asking for the time is accomplished by the sentence: The answer is: Half en kwart. The half hours are indicated differently in Dutch: The quarter hours are similar For more on time telling, have look at the practice lesson 3A. Syntax 3-1 ~ Some more word order: inversion. We have seen that inversion of subject and verb is used to create a question: However, recall from the conversation that inversion happens for other reasons. These are not questions, still there is inversion. The reason is that the adverb "daarna" or the adverbial expression "op een middag" was put before the "subject" + "verb" part for emphasis. This causes inversion. We could also have said: Notice that the verb loses final -t when using the informal second person jij or je in such cases as it does in questions: Grammatica 3-3 ~ Introduction to "naamwoorden". Dutch grammar uses the word "naamwoord" ("lit." name-word) that does not translate well into English. "Naamwoorden" indicates a rather broad class of words, both independently used (like nouns) or used to specify another word (like adjectives). Dutch grammar is therefore structured a bit differently from the English one. Besides "naamwoorden" there are two other large classes of words in Dutch: "werkwoorden" (verbs) and "bijwoorden" (adverbs). A is a fundamental part of speech, occurring in sentences in two different ways: as subjects (performers of action), or objects (recipients of action). As a generality, a noun is the name of a "person, place, thing or concept". Nouns are classified into The latter group is often considered a separate class of words. They stand in for (pro-, voor-) nouns. Words like "hij" - "he" are known as personal pronouns ("persoonlijke voornaamwoorden") Dutch has its own grammatical nomenclature and to use dictionaries and grammars it is useful to know it. Noun is rendered as "zelfstandig naamwoord" ('nameword that stands on itself'). An adjective is called "bijvoeglijk naamwoord" (nameword that can be added). "Naamwoord" is more general than noun. It derives from the Latin term "nomen": nomen substantivum (zelfstandig naamwoord) and nomen adiectivum (bijvoeglijk naamwoord). Adjectives are usually added to nouns to further determine them: Some pronouns, e.g. possessive pronouns ("bezittelijk voornaamwoord") are used as adjectives: A special class of adjectives is formed by the articles ("lidwoorden"): Gender of Nouns. We have seen evidence of word gender in the pronouns we have been encountering; notably 'he', 'she', and 'it' in English and hij, zij, and het in Dutch. We also saw that adjectives depend on gender in Dutch. There are a few rules that help to determine a noun's gender, but mostly it must be learned as children do: word by word. Noun gender is also reflected in the "definite article" It should always be learned as "part of the noun", as this is a good way to memorize gender. Definite Articles. Definite articles are equivalent to an English 'the', and the two basic gender forms in Dutch are as follows: Animate nouns. Much like in English there are three genders for animate nouns (people, pets etc.) and this shows up clearly in their personal pronouns: "hij, zij and het" (he, she and it) and their possessive pronouns "zijn, haar and zijn" (his, her, its): However, "zijn" is not used much anymore to refer to a neuter word and we will see a different way of expressing "its" later. In the plural the gender distinctions are absent: "de mannen, de vrouwen, de kalveren" are all referred to by "zij" (they) and "hun" (their). As you see the definite article is the same for masculine and feminine, but it is not just definite articles, but also adjectives and pronouns that must match the gender of the noun they are related to. Inanimate nouns. In the Netherlands (the North) the distinction between masculine and feminine was lost for "inanimate" nouns (things, concepts etc.) in the 17th century. The feminine and the masculine have merged into a "common" gender north of "de grote riveren" (the Great Rivers: the Meuse, the Rhine and its branches) almost entirely. Someone learning the language therefore best considers Dutch a two-gender language for anything but persons: This does not hold for the South, where a "de klok" may still be referred to as "zij" (she), but it is acceptable standard Dutch to disregard the masculine-feminine distinction. By contrast, the twofold split common-neuter is still very much alive in Dutch and this must be mastered by any beginner to learn the language well. Therefore, it is important when learning Dutch nouns to always learn them together with their correct definite article. That is: This is by far the most important thing you should do right now. The fine distinctions between the varieties of the language can wait. Referring to inanimate common gender. As we saw above the personal pronouns (hij,zij,het) still show the three-gender distinction that Dutch inherited from its Indoeuropean ancestry. That makes it hard to use personal pronouns for an inanimate common gender word. In the South "de klok" may still be called a "she", but Northerners avoid such references and so should you. Strictly speaking it would be correct for Northerners to call a clock a "he", but they often avoid that as well. Nowadays "hij" and "zij" are pretty much restricted to people or their pets, so they indicate "natural" rather than "grammatical" gender, certainly in the North. Notice that Northerners cannot resort to "het" (it) as done in English, because "de klok" is not neuter... This leaves roughly two thirds of all inanimate nouns without a personal pronoun to refer them by. For possessive pronouns (his, her, its) a similar problem exists. We shall see three common ways that speakers use to avoid inanimate hij/zij references: These three aspects of the language have come to play a more prominent role in Dutch than they do in English. One could say that the merger of m/f into common gender has triggered a number of shifts in the language, that for example German or English do not have and must be mastered to speak Dutch well. Rules for gender. There are a few general (and helpful) rules for gender: Another helpful fact is that all genders behave the same in the plural, all use "de", "die", "zij" etc. Apart from these general rules, nouns should be memorized together with their definite article. So, learn "de klok", not just "klok" and "het paard" not just "paard" Double gender. There is an interesting group of words for which the natural gender is in conflict with the grammatical gender, e.g. diminutives of people: Grammatically they are neuter and their articles, adjectives and demonstratives follow the neuter pattern. However the personal and possessive pronouns follow the natural gender: Woordenlijst 3. &lt;br clear="all"&gt; Also included in the vocabulary for Lesson 3 are the ordinal and cardinal numbers 1 through 12 from the table at the beginning of this lesson. Quizlet. You can practice your vocabulary at Quizlet (44 terms) Progress made. If you have studies this lesson well, you should Cumulative term count: 

To make your first encounter with Dutch as easy as possible, the first conversation is given here in several versions: Dutch version. Gesprek 1-1 ~ Vrienden&lt;br&gt; Jan komt Karel op straat tegen. Ze zijn vrienden. Translated into English. Gesprek 1-1 ~ Vrienden&lt;br&gt; Jan meets Karel in the street. They are friends. Glossed version.  Gesprek 1-1 ~ Vrienden  conversation 1-1 ~ Friends  Jan komt Karel op straat tegen. Ze zijn vrienden.  Jan comes Karel on street across. They are friends.  Jan: Goeden-dag, Karel. Hoe gaat het met je?  Jan: Good-day, Karel. How goes it with you?  Karel: Goeden-dag. Dank je, met mij gaat het goed. En met jou?  Karel: Good-day. Thank you, with me goes it good. And with you?  Jan: Dank je, met mij gaat het goed. Tot ziens.  Jan: Thank you, with me goes it good. Until seeing.  Karel: Tot ziens, Jan!  Karel: Until seeing, Jan! Pronunciation: English (more or less). (Note: The G is pronounced loud "from the back of the throat", the same way Spanish-speakers pronounce the "j")&lt;br&gt; Guh-SPREK 1-1 ~ VREEN-duh&lt;br&gt; Yahn comt KAA-rul op straat TAY-Gun. Zuh zian VREEN-duh. 

Vietnamese creates comparatives and superlatives out of adjectives and adverbs simply by adding a word to the end. First of all, what are and ? Comparatives are demonstrated in the following sentences: They simply express a comparison between two items. Superlatives are demonstrated in the following sentences: They express the outcome of a comparison between more than two items. (All of the items being compared, except for "she," are implied in the second sentence.) Comparatives. In Vietnamese, comparatives are formed by placing "hơn" after the adjective or adverb. For example: "(More to come.)" 

Complex analysis is the study of functions of complex variables. Complex analysis is a widely used and powerful tool in certain areas of electrical engineering, and others. Before we begin, you may want to review Complex numbers Complex Numbers. Complex Numbers Complex Functions. A function of a complex variable is a function that can take on complex values, as well as strictly real ones. For example, suppose f(z) = z2. This function sets up a correspondence between the complex number z and its square, z2, just like a function of a real variable, but with complex numbers. Note that, for f(z) = z2, f(z) will be strictly real if z is strictly real. Generally we can write a function f(z) in the form f(z) = f(x+iy) = a(x,y) + ib(x,y), where a and b are real-valued functions. Limits and continuity. As with real-valued functions, we have concepts of limits and continuity with complex-valued functions also – our usual delta-epsilon limit definition: Note that ε and δ are real values. This is implicit in the use of inequalities: only real values are "greater than zero". One difference between this definition of limit and the definition for real-valued functions is the meaning of the absolute value. Here we mean the complex absolute value instead of the real-valued one. Another difference is that of how "z" approaches "w". For real-valued functions, we would only be concerned about "z" approaching "w" from the left, or from the right. In a complex setting, "z" can approach "w" from any direction in the two-dimensional complex plane: along any line passing through "w", along a spiral centered at "w", etc. For example, let formula_7. Suppose we want to show that the formula_8. We can write "z" as formula_9 where we think of "γ" being a small complex quantity. Note then that formula_10. Then, with "L" in our definition being -1, and "w" being i, we have By the triangle inequality, this last expression is less than In order for this to be less than ε, we can require that Thus, for any formula_14, if formula_15, and formula_16, then formula_17. Hence, the limit of formula_7 as "z" approaches i is -1. Differentiation and Holomorphic Functions. Since we have limits defined, we can go ahead to define the "derivative" of a complex function, in the usual way: provided that the limit is the same no matter how Δ"z" approaches zero (since we are working now in the complex plane, we have more freedom!). If such a limit exists for some value "z", or some set of values - a "region", we call the function "holomorphic" at that point or region. Continuity and being single-valued are necessary for being analytic; however, continuity and being single-valued are not sufficient for being analytic. Many elementary functions of complex values have the same derivatives as those for real functions: for example D "z"2 = 2"z". Problem set. Given the above, answer the following questions. Answers. 1. formula_20 2. formula_21 Cauchy-Riemann Equations. We might wonder which sorts of complex functions are in fact differentiable. It would appear that the criterion for holomorphicity is much stricter than that of differentiability for real functions, and this is indeed the case. Suppose we have a complex function where "u" and "v" are real functions. Assume furthermore that "u" and "v" are differentiable functions in the real sense. Then we can let formula_23 in the definition of differentiability approach 0 by varying only "x" or only "y". Therefore "f" can only be differentiable in the complex sense if In fact, if "u" and "v" are differentiable in the real sense and satisfy these two equations, then "f" is holomorphic. These two equations are known as the Cauchy-Riemann equations. Integration. In single variable Calculus, integrals are typically evaluated between two real numbers On the real line, there is one way to get from formula_26 to formula_27. In the complex plane, however, there are infinitely many different paths which can be taken between two points, formula_28 and formula_29. For this reason, complex integration is always done over a path, rather than between two points. Let formula_30 be a path in the complex plane parametrized by formula_31, and let formula_32 be a complex-valued function. Then the contour integral is defined analogously to the line integral from multivariable calculus: Example Let formula_34, and let formula_30 be a line from 0 to 1+i. This curve can be parametrized by formula_36, with formula_37 ranging from 0 to 1. Now we can compute Note that we also have This indicates that complex antiderivatives can be used to simplify the evaluation of integrals, just as real antiderivatives are used to evaluate real integrals. Cauchy's Theorem. Cauchy's theorem states that if a function formula_40 is holomorphic in the closure of an open set formula_41, and formula_30 is a simple closed curve in formula_41, then This can be understood in terms of Green's theorem, though this does not readily lead to a proof, since Green's theorem only applies under the assumption that f has continuous first partial derivatives... Contour Integration. Cauchy's theorem allows for the evaluation of many improper real integrals (improper here means that one of the limits of integration is infinite). As an example, consider Since we consider We now integrate over the indented semicircle contour, pictured above. We parametrize each segment of the contour as follows By Cauchy's Theorem, the integral over the whole contour is zero. So, We now handle each of these integrals separately. Recalling the definition of the sine of a complex number, Now we evaluate the other two integrals As formula_61, the integrand approaches one, so The fourth integral is equal to zero, but this is somewhat more difficult to show. Its form is similar to that of the third segment: This integrand is more difficult, since it need not approach zero everywhere. This difficulty can be overcome by splitting up the integral, but here we simply assume it to be zero. Combining everything, we now have Hence, Cauchy's Integral Formula. Cauchy's integral formula characterizes the behavior of holomorphics functions on a set based on their behavior on the boundary of that set. If formula_41 is an open set with a piecewise smooth boundary and formula_40 is holomorphic in formula_68, then This is a remarkable fact which has no counterpart in multivariable calculus. It says that if we know the values of a holomorphic function along a closed curve, then we know its values everywhere in the interior of the curve. Because formula_70, an open set, it follows that formula_71 for all formula_72. Hence the integrand in Cauchy's integral formula is infinitely differentiable with respect to z, and by repeatedly taking derivatives of both sides, we get This result shows that holomorphicity is a much stronger requirement than differentiability. In the complex plane, if a function has just a single derivative in an open set, then it has infinitely many derivatives in that set. Corollaries of Cauchy's Theorem. Cauchy's Theorem and integral formula have a number of powerful corollaries: 

Alkynes are hydrogenated in generally the same way as alkenes. However, standard catalysts like Pd/C will not allow hydrogenation to stop before the alkane stage. A "poisoned catalyst" will permit the reduction of a triple bond to a double bond by "syn" addition, but no further. A common poisoned catalyst is Lindlar's catalyst: Pd/CaCO3 treated with quinoline and lead acetate. "Anti" addition of hydrogen atoms can be achieved with a dissolving metal reduction. Dissolving Metal Reduction. The dissolving metal reaction takes place in a solution of sodium or lithium metal dissolved in liquid ammonia. An alkali metal is dissolved in liquid ammonia and forms a solution containing solvated electrons. Once an alkyne is added to the solution, the electrons add to the antibonding π molecular orbital and produce a radical-anion intermediate. This intermediate has one unpaired electron and has a negative charge. The radical anion then reacts with a proton source (generally ethanol or tertiary-butyl alcohol which is added to the ammonia) to yield a vinyl radical. The radical then accepts another electron and forms a carbanion. The alkene is formed when the vinyl carbanion is protonated. Protonation normally occurs from the side of the double bond that is least sterically hindered, so a trans-alkene is produced. 

There are several ways to halogenate Alkenes. The simplest is by adding appropriate hydrogen halides, where the reaction follows Markovnikov's rule. This rule states that the halogen will become bonded to the more highly substituted carbon. This is a result of greater carbocation stability of more highly substituted carbon atoms, which have more ability to distribute the positive charge to the attached alkyl groups. When both carbon atoms have the same number of substituents, a mixture of products will often result. 

Haskell is a functional programming language. It is distinct in a few ways: Haskell is enjoyable to use because dealing with pure functions makes code much easier to reason about, and the advanced type system helps catch silly and profound mistakes. Our aim in this book is to introduce you to the Haskell programming language — from the very basics to advanced features — and to computer programming in general. We urge seasoned programmers to be especially patient with this process. The languages you are familiar with are likely to differ greatly from Haskell, and the habits acquired from those languages might make it difficult to understand how things work — Haskell is simple, but different. Learning to see the world through the warped mindset of a functional programmer is an adventure in a brave new world, which brings knowledge valuable far beyond the boundaries of any language. Overview. The book is divided into three sections: a Beginner's Track, an Advanced Track, and a section called "Haskell in Practice". The last section, which covers practical use cases, depends only on the Beginner's Track. Seasoned programmers may read the overview to quickly evaluate what makes Haskell unique and different from other languages. Beginner's Track. This section introduces you to Haskell basics and some commonly used libraries. At the end of this track you should be able to write simple Haskell programs. Most chapters include exercises, with solutions, for your practice. Advanced Track. This section introduces wider functional programming concepts such as different data structures and type theory. It will also cover more practical topics like concurrency. Haskell in Practice. Day-to-day issues of working with Haskell include items such as knowing the standard libraries, building graphical interfaces, and working with databases. You should be able to jump directly to this section from the beginner's track. Tutorials that have been incorporated into the Haskell Wikibook. The following may be read independently, but their content has been imported and adapted already into the Wikibook here 



= Lesson 7: 这是什么? What's this? = Text 1. "You can check out the translations here." Vocabulary. Stroke orders. "More will be added if it's helpful." Grammar. Chinese Names. &lt;br&gt; In this lesson, we learn how to say "something is something" in Chinese. The first thing you need to know is that the sentence structure of Chinese is very similar to that of English in that they both follow the pattern of Subject-Verb-Object (SVO). But unlike many Western languages, verbs in Chinese aren't conjugated and noun and adjective endings don't change. They are never affected by things such as time or person. 这(/那)是什么？. This sentence means "What's this/that?": The sentences, if broken down literally, shows that the ordering of words differs in English and Chinese: The order of the sentences may seem a little bit tricky, but don't worry about that, we will discuss this later. A　是　B. This sentence means "A is B." "是" (shì), the equational verb "to be", can be used as the English "is" or "equals". When used in a simple Subject-Verb-Object sentence, the subject defines the object. Since Chinese verbs never change, no other forms for shì exist such as "was" or "am" in English. Also, articles like "a" and "the" are absent in Chinese. They are not translated. For example: A　不是　B. This sentence means "A is not B." in which shì is negated when preceded by "不" (bu). "不" literally means "no", "not". For example: Now, we come back to the "what's this/that?" questions. We already see that the order is a bit tricky comparing to the English question order. But comparing to the latter pattern "A　是　B", we find the similarity: their orders are identically the same. In fact, like particles, question words make statements into questions without changing the order of the sentence. To make one, simply substitute the QW in where the subject would be in the answer. Comparison: 吗. "吗"(ma) is a final interrogative particle used to form a question sentence. Adding this character at the end of a statement transforms the sentence into a question. Example 1: Example 2: 是/不. "是" (shì) can be used to answer a simple yes/no question. In this case, "是" means "yes", whilst "不" (bú) or "不是" (bú shì) means "no" (literally, "not is"). How to answer yes/no questions correctly in Chinese? Usually, it's the same as in English, but pay attention if the questions are negative, like "Isn't this a book?". In Chinese, you answer to the questions, not the fact. If the question itself is a negative answer, use "不是" or simply "不", vice versa. For example: A asks if that's a book in a negative way. If the object is not a book, you should nevertheless approve A's saying first. So we use "是" to acknowledge that A is correct, and then say "this is not (a) book" to emphasis A is right; In the case of that is a book, you should deny A's saying first, using "不" (no) to point out A is wrong, then make a new statement by noting that "这是书" (this "is" (a) book). One more example: 的. Character "的" (de) indicates that the previous word has possession of the next one. In English it functions like "'s" or like the word "of" but with the position of possessor and possessee switched. For example: Further reading. Read the following article, and then answer the questions in Chinese. Questions: 

Energy. Energy as described on Wikipedia is "the property that must be transferred to an object in order to perform work on, or to heat, the object". Energy is a conserved quantity, the Law of Conservation of Energy on Wikipedia states that "energy can be converted in form, but not created or destroyed". Common forms of energy in physics are potential and kinetic energy. The potential energy is usually the energy due to matter having certain position (configuration) in a field, commonly the gravitational field of Earth. Kinetic energy is the energy due to motion relative to a frame of reference. In thermodynamics, we deal with mainly work and heat, which are different manifestations of the energy in the universe. Work. Work is said to be done by a system if the effect on the surroundings can be reduced solely to that of lifting a weight. Work is only ever done at the boundary of a system. Again, we use the intuitive definition of work, and this will be complete only with the statement of the second law of thermodynamics. Consider a piston-cylinder arrangement as found in automobile engines. When the gas in the cylinder expands, pushing the piston outwards, it does work on the surroundings. In this case work done is mechanical. But how about other forms of energy like heat? The answer is that heat cannot be completely converted into work, with no other change, due to the second law of thermodynamics. In the case of the piston-cylinder system, the work done during a cycle is given by "W", where "W = −∫ F dx = −∫ p dV", where "F = p A", and "p" is the pressure on the inside of the piston (note the minus sign in this relationship). In other words, the work done is the area under the "p-V" diagram. Here, "F" is the external opposing force, which is equal and opposite to that exerted by the system. A corollary of the above statement is that a system undergoing free expansion does no work. The above definition of work will only hold for the quasi-static case, when the work done is reversible work. A consequence of the above statement is that work done is not a state function, since it depends on the path (which curve you consider for integration from state 1 to 2). For a system in a cycle which has states 1 and 2, the work done depends on the path taken during the cycle. If, in the cycle, the movement from 1 to 2 is along "A" and the return is along "C", then the work done is the lightly shaded area. However, if the system returns to 1 via the path "B", then the work done is larger, and is equal to the sum of the two areas. The above image shows a typical indicator diagram as output by an automobile engine. The shaded region is proportional to the work done by the engine, and the volume "V" in the "x"-axis is obtained from the piston displacement, while the "y"-axis is from the pressure inside the cylinder. The work done in a cycle is given by "W", where formula_1 Work done by the system is negative, and work done on the system is positive, by the convention used in this book. Flow Energy. So far we have looked at the work done to compress fluid in a system. Suppose we have to introduce some amount of fluid into the system at a pressure "p". Remember from the definition of the system that matter can enter or leave an open system. Consider a small amount of fluid of mass "dm" with volume "dV" entering the system. Suppose the area of cross section at the entrance is "A". Then the distance the force "pA" has to push is "dx = dV/A". Thus, the work done to introduce a small amount of fluid is given by "pdV", and the work done per unit mass is "pv", where "v = dV/dm" is the specific volume. This value of "pv" is called the "flow energy". Examples of Work. The amount of work done in a process depends on the irreversibilities present. A complete discussion of the irreversibilities is only possible after the discussion of the second law. The equations given above will give the values of work for quasi-static processes, and many real world processes can be approximated by this process. However, note that work is only done if there is an opposing force in the boundary, and that a volume change is not strictly required. Work in a Polytropic Process. Consider a polytropic process "pVn=C", where "C" is a constant. If the system changes its states from 1 to 2, the work done is given by formula_2 And additionally, if n=1 formula_3 Heat. Before thermodynamics was an established science, the popular theory was that heat was a fluid, called "caloric", that was stored in a body. Thus, it was thought that a hot body transferred heat to a cold body by transferring some of this fluid to it. However, this was soon disproved by showing that heat was generated when drilling bores of guns, where both the drill and the barrel were initially cold. Heat is the energy exchanged due to a temperature difference. As with work, heat is defined at the boundary of a system and is a path function. Heat rejected by the system is negative, while the heat absorbed by the system is positive. Specific Heat. The "specific heat" of a substance is the amount of heat required for a unit rise in the temperature in a unit mass of the material. If this quantity is to be of any use, the amount of heat transferred should be a linear function of temperature. This is certainly true for ideal gases. This is also true for many metals and also for real gases under certain conditions. In general, we can only talk about the average specific heat, "cav = Q/mΔT". Since it was customary to give the specific heat as a property in describing a material, methods of analysis came to rely on it for routine calculations. However, since it is only constant for some materials, older calculations became very convoluted for newer materials. For instance, for finding the amount of heat transferred, it would have been simple to give a chart of "Q(ΔT)" for that material. However, following convention, the tables of "cav(ΔT)" were given, so that a double iterative solution over "cav" and "T" was required. Calculating specific heat requires us to specify what we do with Volume and Pressure when we change temperature. When Volume is fixed, it is called specific heat at constant volume (Cv). When Pressure is fixed, it is called specific heat at constant pressure (Cp). Latent Heat. It can be seen that the specific heat as defined above will be infinitely large for a phase change, where heat is transferred without any change in temperature. Thus, it is much more useful to define a quantity called "latent heat", which is the amount of energy required to change the phase of a unit mass of a substance at the phase change temperature. Adiabatic Process. An "adiabatic" process is defined as one in which there is no heat transfer with the surroundings, that is, the change in amount of energy dQ=0 A gas contained in an insulated vessel undergoes an adiabatic process. Adiabatic processes also take place even if the vessel is not insulated if the process is fast enough that there is not enough time for heat to escape ("e.g." the transmission of sound through air). Adiabatic processes are also ideal approximations for many real processes, like expansion of a vapor in a turbine, where the heat loss is much smaller than the work done. First Law of Thermodynamics. Joule Experiments. It is well known that heat and work both change the energy of a system. Joule conducted a series of experiments which showed the relationship between heat and work in a thermodynamic cycle for a system. He used a paddle to stir an insulated vessel filled with fluid. The amount of work done on the paddle was noted (the work was done by lowering a weight, so that work done = "mgz"). Later, this vessel was placed in a bath and cooled. The energy involved in increasing the temperature of the bath was shown to be equal to that supplied by the lowered weight. Joule also performed experiments where electrical work was converted to heat using a coil and obtained the same result. Statement of the First Law for a Closed System. The first law states that "when heat and work interactions take place between a closed system and the environment, the algebraic sum of the heat and work interactions for a cycle is zero". Mathematically, this is equivalent to  dQ + dW = 0 for any cycle closed to mass flow Q is the heat transferred, and W is the work done on or by the system. Since these are the only ways energy can be transferred, this implies that the total energy of the system in the cycle is a constant. One consequence of the statement is that the total energy of the system is a property of the system. This leads us to the concept of internal energy. Internal Energy. In thermodynamics, the "internal energy" is the energy of a system due to its temperature. The statement of first law refers to thermodynamic cycles. Using the concept of internal energy it is possible to state the first law for a non-cyclic process. Since the first law is another way of stating the conservation of energy, the energy of the system is the sum of the heat and work input, "i.e.", "ΔE = Q + W". Here "E" represents the internal energy (U) of the system along with the kinetic energy (KE) and the potential energy (PE) and is called the "total energy" of the system. This is the statement of the first law for non-cyclic processes, as long as they are still closed to the flow of mass ("E = U + KE + PE"). The KE and PE terms are relative to an external reference point i.e. the system is the gas within a ball, the ball travels in a trajectory that varies in height H and velocity V and subsequently KE and PE with time, but this has no affect upon the energy of the gas molecules within the ball, which is dictated only by the internal energy of the system (U). Thermodynamics does not define the nature of the internal energy, but it can be rationalised using other theories (i.e. the gas kinetic theory), but in this case is due to the KE and PE of the gas molecules within the ball, not to be mistaken with the KE and PE of the ball itself. For gases, the value of "KE" and "PE" is quite small, so the important term is the internal energy function "U". In particular, since for an ideal gas the state can be specified using two variables, the state variable "u" is given by "u(v, T)", where "v" is the specific volume and "T" is the temperature. Introducing this temperature dependence explicitly is important in many calculations. For this purpose, the constant-volume heat capacity is defined as follows: "cv = (∂u/∂t)v", where "cv" is the specific heat at constant volume. A constant-pressure heat capacity will be defined later, and it is important to keep them straight. The important point here is that the other variable that U depends on "naturally" is v, so to isolate the temperature dependence of U you want to take the derivative at constant v. Internal Energy for an Ideal Gas. In the previous section, the internal energy of an ideal gas was shown to be a function of both the volume and temperature. Joule performed an experiment where a gas at high pressure inside a bath at the same temperature was allowed to expand into a larger volume. In the above image, two vessels, labeled A and B, are immersed in an insulated tank containing water. A thermometer is used to measure the temperature of the water in the tank. The two vessels A and B are connected by a tube, the flow through which is controlled by a stop. Initially, A contains gas at high pressure, while B is nearly empty. The stop is removed so that the vessels are connected and the final temperature of the bath is noted. The temperature of the bath was unchanged at the end of the process, showing that the internal energy of an ideal gas was the function of temperature alone. Thus "Joule's law" is stated as "(∂u/∂v)t = 0". Enthalpy. According to the first law, "dQ + dW = dE" If all the work is pressure volume work, then we have "dW = − p dV" ⇒ "dQ = dU + pdV = d(U + pV) - Vdp" ⇒ "d(U + pV) = dQ + Vdp" We define "H ≡ U + pV" as the "enthalpy" of the system, and "h = u + pv" is the specific enthalpy. In particular, for a constant pressure process, "ΔQ = ΔH" With arguments similar to that for "cv", "cp = (∂h/∂t)p". Since "h", "p", and "t" are state variables, "cp" is a state variable. As a corollary, for ideal gases, "cp = cv + R", and for incompressible fluids, "cp = cv" Throttling. "Throttling" is the process in which a fluid passing through a restriction loses pressure. It usually occurs when fluid passes through small orifices like porous plugs. The original throttling experiments were conducted by Joule and Thompson. As seen in the previous section, in adiabatic throttling the enthalpy is constant. What is significant is that for ideal gases, the enthalpy depends only on temperature, so that there is no temperature change, as there is no work done or heat supplied. However, for real gases, below a certain temperature, called the "inversion point", the temperature drops with a drop in pressure, so that throttling causes cooling, "i.e.", "p1 &lt; p2 ⇒ T1 &lt; T2." The amount of cooling produced is quantified by the Joule-Thomson coefficient "μJT = (∂T/∂p)h". For instance, the inversion temperature for air is about 400 °C. 

Leaves (forms) Chapter 3. Plant Cells and Tissues&lt;br&gt;Leaves (forms) Leaf Shapes, Margins, Apices, and Bases. Simple, Angiosperm Leaves. This page presents the terminology used to describe the shapes and forms of leaves. The following examples are all simple leaves. All are from flowering plants (angiosperms). Descriptive terms are given for the lamina shape (that is, its overall shape in outline), the lamina margin, the form or shape of both the apex (tip) and the base (lamina junction with petiole) of the leaf, and (in some cases) other qualities evident in the example picture presented. Compound, Angiosperm Leaves. These examples are all compound leaves. In these leaves, it is not always very instructive to describe the "leaf outline" and better to describe the number (or range) of leaflets. If the number is odd, there will be an unpaired leaflet at the tip. The shape terminology can be applied to the individual leaflets. &lt;BR&gt; &lt;HR&gt; Links. Botanical Terms 

__NOEDITSECTION__ Presentation. UK Constitution and Government 

William I. The course of English political, legal and cultural history was changed in 1066, when William, Duke of Normandy (also called William the Conqueror) successfully invaded the nation and displaced the Saxon king, Harold II. In 1066 King Edward, also called St Edward the Confessor, died. His cousin, the Duke of Normandy, claimed that the childless King had named him heir during a visit to France, and that the other claimant to the throne, Harold Godwinson, had pledged to support William when he was shipwrecked in Normandy. The veracity of this tale, however, is doubtful, and Harold took the crown upon King Edward's death. William, however, invaded England in September, and defeated (and killed) Harold at the famous Battle of Hastings in October. William II. In 1087, King William I died, and divided his lands and riches between his three sons. The eldest, Robert, became Duke of Normandy; the second, William, became King of England; the youngest, Henry, received silver. Henry, however, eventually came to possess all of his father's dominions. William II died without children, so Henry became King. Henry later invaded Normandy, imprisoned his brother, and took over the Duchy of Normandy. Henry I, Stephen and Matilda. Henry, whose sons had predeceased him, took an unprecedented step: naming a woman as his heir. He declared that his daughter Matilda would be the next Queen. However, Matilda's claim was disputed by Stephen, a grandson of William I in the female line. After Henry I died in 1135, Stephen usurped the throne, but he was defeated and imprisoned by Matilda in 1141. Later, however, Matilda was defeated, and Stephen took the throne. Matilda, however, was not completely defeated. She escaped from Stephen's army, and her own son, Henry Plantagenet, led a military expedition against Stephen. Stephen was forced to agree to name Henry as his heir, and when Stephen died in 1154, Henry took the throne, commencing the Plantagenet dynasty. 

Henry II. With the death of King Stephen, Henry Plantagenet took the throne as King Henry II. He already had control over the duchy of Normandy; he had also inherited Anjou from his father Geoffrey. Furthermore, he acquired many territories from his wife, Eleanor of Aquitaine. Henry thus had a vast territory when he came to the throne; as King of England, he took over Ireland. Henry II made other remarkable achievements in England. He established courts throughout England and introduced trial by jury. Furthermore, he reduced the power of ecclesiastical courts. The Archbishop of Canterbury and Lord High Chancellor, Thomas à Becket, opposed the King's attempt to take power from the Church. At a confrontation between the two in 1170, Henry II famously said, "Who will rid me of this turbulent priest?" Four of his knights took him literally, and in December murdered Becket. Henry, however, did not have good relations with his sons. In 1170, his eldest son Henry was crowned, and is known as Henry the Young King. In 1173, the Young King and his brothers revolted against Henry II, planning to dethrone him and leave the Young King as the sole ruler in England. In 1174, the revolt failed, and all of the brothers surrendered. Later, in 1189, Henry II's third son, Richard, attacked and defeated him. Henry II died days after his defeat, and Richard, nicknamed "the Lionheart," became King. Richard I. Richard the Lionheart is often portrayed as a hero, but he did not do much for England. In fact, he spent almost all of his time outside the nation, and did not even find it necessary to learn English. He is most famous for his fighting in the Crusades, a holy war seeking to assert Christian dominance over Jerusalem. John. Richard's successor was his brother, John Allin. Henry II had granted John the lands of Ireland, so when John came to the throne, the titles Lord of Ireland and King of England were united. However, though Ireland became a dominion of the Crown, several lands on the Continent, including most of Normandy, were lost during John's reign. King John was very unpopular with the nation's magnates, the barons, whom he taxed. A particularly resented tax was the scutage, a penalty paid by barons who failed to supply the King with military resources. In 1215, after John had been defeated in France, several barons rebelled. Later in that year, John compromised and signed the "Magna Carta", or Great Charter. It guaranteed political liberties and provided for a church free from domination by the monarchy. These liberties and privileges, however, were not extended to the common man; rather, they were granted to the barons. Nonetheless, the document is immensely significant in English constitutional history as it is a major indication of a limitation on the power of the Crown. King John, however, broke the provisions of the Charter later, claiming that he agreed to it under duress. In the next year, when he was retreating from a French invasion, John lost England's most valuable treasures - the Crown Jewels - in a marsh known as The Wash. His mental and physical health deteriorated, and he later died from dysentery. Henry III. John was succeeded by his son, Henry, who was only nine years old. Henry III, despite a reign that lasted over half a century, is not a particularly memorable or noteworthy monarch. Nonetheless, a very significant political development occurred during Henry III's reign. In 1258, one of Henry's opponents, Simon de Montfort, called a Parliament, the forerunner of the modern institution. It, however, bears little resemblance to the modern body, as it had little power. Simon de Montfort, who was married to Henry III's sister, defeated and imprisoned his brother-in-law in 1264. He was originally supported by Henry's son Edward, but the latter later returned to his father's side. Edward defeated de Montfort in 1265 at the Battle of Evesham and restored Henry III. In 1270, the ageing Henry gave up most of power to his son; two years later, he died, and Edward succeeded to the throne. Edward I. Edward I was the monarch who brought the entire British Isles under English domination. In order to raise money in the war against the rebellious Wales, Edward instituted a tax on Jewish moneylenders. The tax, however, was too high for the moneylenders, who eventually became too poor to pay. Edward accused them of disloyalty and abolished the right of Jews to lend money. He also ordered that all Jews wear a yellow star on their clothing; that idea was later adopted by Adolf Hitler in Germany. Edward also executed hundreds of Jews, and in 1290 banished all of them from England. In 1291, the Scottish nobility agreed to submit to Edward. When Queen Margaret I died, the nobles allowed Edward to choose between the rival claimants to the throne. Edward installed the weak John Balliol as monarch, and easily dominated Scotland. The Scots, however, rebelled. Edward I executed the chief dissenter, William Wallace, further antagonising Scotland. Edward II. When Edward I died in 1307, his son Edward became King. Edward II abandoned his father's ambitions to conquer Scotland. Furthermore, he recalled several men his father had banished. The barons, however, rebelled against Edward. In 1312, Edward agreed to hand over power to a committee of barons known as "ordainers." These ordainers removed the power of representatives of commoners to advise the monarch on new laws, and concentrated all power in the nobility. Meanwhile, Robert the Bruce was slowly reconquering Scotland. In 1314, Robert's forces defeated England's in battle, and Robert gained control over most of Scotland. In 1321, the ordainers banished a baron allied with the King, Hugh le Despencer, along with his son. In 1322, Edward reacted by recalling them and attacking the barons. He executed the leader of the ordainers, the Earl of Lancaster, and permitted the Despencers to rule England. The Despencers declared that all statutes created by the ordainers were invalid, and that thereafter, no law would be valid unless it had received the assent of the Commons, representatives of the commoners of England. However, the Despencers became corrupt, causing them to be very unpopular, even with Edward's own wife, Isabella. In 1325, Isabella went to France, and in 1326, she returned, allied with Roger Mortimer, one of the barons Edward had defeated. The two killed the Despencers and forced Edward to resign his crown to his son, also named Edward. Edward II was imprisoned and later killed. Edward III. Since Edward III was a child, Isabella and Roger Mortimer ruled England in his stead. When Edward III became eighteen, however, he had Mortimer executed and banished his mother from court. In 1328, when Charles IV, Isabella's father and King of France, died, Edward claimed France, suggesting that the kingdom should pass to him through his mother. His claim was opposed by Philip VI, who claimed that the throne could only pass in the male line. Edward declared war on Philip, setting off the Hundred Years' War. The British claim to the French throne was not abandoned until the nineteenth century. Richard II. Richard II succeeded his grandfather, Edward III, in 1377. Richard II was only about ten years old when coming to the throne. Even as an adult, Richard II was a rather weak king. In 1399, he was deposed by his cousin, Henry of Bolingbroke, and probably murdered the next year. 



Editing Chapters. You simply click on the "edit this page" tab at the top of any page you want to edit. You can also just edit a section or subsection - any place you see a little 'edit' to the far right above a line, click on it to edit that section. The thoughts behind the editing process of the chapters will be documented in the "discussion" tab of each chapter - go there to record your intentions for your assigned chapter. If you would like to make any suggestions to the editorial board, we welcome them. At the top of the XML: Managing Data Exchange contents page is a tab to "discussion" follow that link and edit the page with your comments, be sure to follow the directions for entries. Consistency guidelines.  - developed to create a consistent look and feel to the book Layout. The bottom of each chapter should include a chapter summary highlighting the key points of the chapter. Follow the formatting guideline detailed later on this page. As new concepts and terminology are introduced in each chapter, make sure that they are defined first before using them; so that the reader understands what the words and ideas represent. Bold the first occurrence of the word (or preceding the definition). Capitalize the first letter of a title and allow the rest of the letters to be lower case, as in this title. For consistency, label all tables, figures (e.g. charts and photos), and code examples as "exhibits." Place the label for exhibits at the top of the exhibit, flush left, and enumerate the exhibits, followed by a brief caption. Examples: "Exhibit 1: XML data types" and "Exhibit 2: Schema code example" Think of your chapter as a data model or a large steak. A useful data model is broken down into its smallest attributes that require recording and maintaining for future analysis. A large steak is best eaten in bite size pieces to prevent choking and promote good digestion. As you revise your chapter consider the basic elements of your topic; break your topic down and present each element in a section - so that it is easy for the reader to follow and understand each part. Feel free to use bullets to represent lists. To continue with the steak analogy, it is also good to eat slowly, taking breaths between each bite. Allow for spaces between each subsection and greater spaces between each section. This opens up the text so that it does not appear to dense and compact - and therefore intimidating. Scroll down the contents page to "Appendices". You will see a link to the chapter "Using an XML editor". When you open this link - go to the edit this page tab and insert a link under the NetBeans heading with the name of your chapter. This is the format for the link: After you save your edit, the XML Editor page will contain a link to your new NetBeans page. It will appear red because no information has been posted to that page yet. Click on the link and paste all the NetBeans information from your chapter into the edit box. In your assigned chapter, at the places where there were NetBeans information (that you just cut and paste into the XML editor page) - put a link to the XML editor - your Chapter page. This help page also has a link to a "sandbox." This is a blank page that allows you to experiment with wiki without affecting a real page. Code examples. If you run into trouble while inserting code, one of the two tags mentioned above might be able to get you out. Exercises &amp; Answers. Create two links at the bottom of your chapter, one to an Exercises page and one to the Answers page already made (you can find this answer page link at the bottom of the contents page). In the Exercises page, cut and paste your chapter exercises into a new page. An easy way to create a new page is simply to go to your main chapter page, type "exercises" at the end of the URL, and hit enter. If the page does not exist, Wikipedia will ask you to create the page by going to "edit this page." At the top and bottom of the Exercises page, provide a link to the chapter and the Answers page—and vice versa with the Answers page. Elements of Style - Principles of Composition. taken from "the Elements of Style" - by William Strunk and E.B. White Source code. For inserting source code (e.g., XML or Java) into the book, use the following format: Use the following wiki code: Section summary. Section summaries should appear in the following format: Use the following wiki code: Story. Stories/case studies should appear in the following format: Use the following wiki code: 

This series is represented by the months. If the month has 31 days, it gets a one. If not, it gets a zero. Alternatively, it can be described as the remainder after dividing by 2 of the length of each month in a common year (as opposed to a leap year). 

Before we begin: This chapter assumes knowledge of Some Counting Problems. "..more to come" Generating functions. "..some motivation to be written" To understand this section you need to see why this is true: For a more detailed discussion of the above, head to Infinity and infinite processes. Generating functions, otherwise known as Formal Power Series, are useful for solving problems like: where how many unique solutions are there if formula_4? Before we tackle that problem, let's consider the infinite polynomial: We want to obtain a "closed form" of this infinite polynomial. The "closed form" is simply a way of expressing the polynomial so that it involves only a finite number of operations. To find the "closed form " we starting with our function: formula_6 We multiply both sides of the function by x to get: formula_7 Next we subtract S-xS to get S - xS = 1 + x + x^2 + x^3 ... - x - x^2 - x^3 ... Grouping like terms we get Which simplifies to (1 - x)S = 1 Dividing both sides by formula_8 we get formula_9 So the closed form of is For convenience we can write, although this isn't true for any particular value of "x". info - Infinite sums. The two expressions are not "equal". It's just that for certain values of "x" (-1 &lt; x &lt; 1), we can approximate the right hand side "as closely as possible" by adding up a large number of terms on the left hand side. For example, suppose "x" = 1/2, RHS = 2; we approximate the LHS using only 5 terms we get LHS equals 1 + 1/2 + 1/4 + 1/8 + 1/16 = 1.9375, which is close to 2, as you can imagine by adding more and more terms, we will get closer and closer to 2. Anyway we really only care about its nice algebraic properties, not its numerical value. From now on we will omit the condition for equality to be true when writing out generating functions. Consider a more general case: where "A" and "B" are constants. We can derive the "closed-form" as follows: The following identity as derived above is worth investing time and effort memorising. Exercises. 1. Find the closed-form: 2. Given the closed-form, find a function f(n) for the coefficients of xn: Method of Substitution. We are given that: and we can obtain many other generating functions by substitution. For example: letting z = x2 we have: Similarly is obtained by letting z = Bx then multiplying the whole expression by A. Exercises. 1. What are the coefficients of the powers of x: 2. What are the coefficients of the powers of x (Hint: take out a factor of 1/2): Linear Recurrence Relations. The Fibonacci series where each and every number, except the first two, is the sum of the two preceding numbers. We say the numbers are "related" if the value a number takes depends on the values that come before it in the sequence. The Fibonacci sequence is an example of a recurrence relation, it is expressed as: where xn is the ("n"+ 1)th number in the sequence. Note that the first number in the sequence is denoted x0. Given this recurrence relation, the question we want to ask is "can we find a formula for the ("n"+1)th number in the sequence?". The answer is yes, but before we come to that, let's look at some examples. Example 1. The expressions define a recurrence relation. The sequence is: 1, 1, 5, 13, 41, 121, 365... Find a formula for the ("n"+1)th number in the sequence. Solution Let G(z) be generating function of the sequence, meaning the coefficient of each power (in ascending order) is the corresponding number in the sequence. So the generating functions looks like this Now, by a series of algebraic manipulations, we can find the closed form of the generating function and from that the formula for each coefficient by definition xn - 2xn - 1 - 3xn - 2 = 0 by the method of partial fractions we get: each part of the sum is in a recognisable closed-form. We can conclude that: the reader can easily check the accuracy of the formula. Example 2. Find a non-recurrent formula for xn. Solution Let G(z) be the generating function of the sequence described above. Therefore xn = 1, for all n. Example 3. A linear recurrence relation is defined by: Find the general formula for xn. Solution Let G(z) be the generating function of the recurrence relation. Therefore Exercises. 1. Derive the formula for the ("n+1")th numbers in the sequence defined by the linear recurrence relations: 2. Derive the formula for the ("n+1")th numbers in the sequence defined by the linear recurrence relations: 3. (Optional) Derive the formula for the ("n+1")th Fibonacci numbers. Further Counting. Consider the equation For a fixed positive integer n, how many solutions are there? We can count the number of solutions: As you can see there are (n + 1) solutions. Another way to solve the problem is to consider the generating function Let H(z) = G(z)G(z), i.e. I claim that the coefficient of zn in H(z) is the number of solutions to a + b = n, a, b &gt; 0. The reason why lies in the fact that "when multiplying like terms, indices add". Consider Let it follows Now the coefficient of zn (for n ≥ 0) is clearly the number of solutions to a + b = n (a, b &gt; 0). We are ready now to derive a very important result: let tk be the number solutions to a + b = n (a, b &gt; 0). Then the generating function for the sequence tk is i.e. Counting Solutions to a1 + a2 + ... + am = n. Consider the number of solutions to the following equation: where ai ≥ 0; i = 1, 2, ... m. By applying the method discussed previously. If tk is the number of solutions to the above equation when n = k. The generating function for tk is but what is tk? Unless you have learnt calculus, it's hard to derive a formula just by looking at the equation of T(z). Without assuming knowledge of calculus, we consider the following counting problem. "You have three sisters, and you have n (n ≥ 3) lollies. You decide to give each of your sisters at least one lolly. In how many ways can this be done?" One way to solve the problem is to put all the lollies on the table in a straight line. Since there are "n" lollies there are ("n" - 1) gaps between them (just as you have 5 fingers on each hand and 4 gaps between them). Now get 2 dividers, from the ("n" - 1) gaps available, "choose" 2 and put a divider in each of the gaps you have chosen! There you have it, you have divided the "n" lollies into three parts, one for each sister. There are formula_50 ways to do it! If you have 4 sisters, then there are formula_51 ways to do it. If you have m sisters there are formula_52 ways to do it. Now consider: "You have three sisters, and you have n lollies. You decide to give each of your sisters some lollies (with no restriction as to how much you give to each sister). In how many ways can this be done?" Please note that you are just solving: where ai ≥ 0; i = 1, 2, 3. You can solve the problem by putting n + 3 lollies on the table in a straight line. Get two dividers and "choose" 2 gaps from the n + 2 gaps available. Now that you have divided n + 3 lollies into 3 parts, with each part having 1 or more lollies. Now take back 1 lollies from each part, and you have solved the problem! So the number of solutions is formula_53. More generally, if you have m sisters and n lollies the number of ways to share the lollies is Now to the important result, as discussed above the number of solutions to where ai ≥ 0; i = 1, 2, 3 ... m is i.e. Example 1. The closed form of a generating function T(z) is and tk in the coefficient of zk is T(z). Find an explicit formula for tk. Solution Therefore tk = k Example 2. Find the number of solutions to: for all positive integers n (including zero) with the restriction a, b, c ,d ≥ 0. Solution By the formula so More Counting. We turn to a slightly harder problem of the same kind. Suppose we are to count the number of solutions to: for some integer formula_61, with "a", "b", also "c" greater than or equal zero. We can write down the closed form straight away, we note the coefficient of "x""n" of: is the required solution. This is due to, again, the fact that when multiplying powers, indices add. To obtain the number of solutions, we break the expression into recognisable closed-forms by method of partial fraction. Example 1. Let sk be the number of solutions to the following equation: Find the generating function for sk, then find an explicit formula for sn in terms of "n". Solution Let T(z) be the generating functions of tk It's not hard to see that Example 2. Let tk be the number of solutions to the following equation: Find the generating function for tk, then find an explicit formula for tn in terms of "n". Solution Let T(z) be the generating functions of tk Exercises. 1. Let be the generating functions for tk (k = 0, 1, 2 ...). Find an explicit formula for tk in terms of "k". 2. How many solutions are there the following equations if "m" is a given constant where a, b and c ≥ 0 Problem Set. 1. A new Company has borrowed $250,000 initial capital. The monthly interest is 3%. The company plans to repay $"x" before the end of each month. Interest is added to the debt on the last day of the month (compounded monthly). Let Dn be the remaining debt after "n" months. a) Define Dn recursively. b) Find the minimum values of "x". c) Find out the general formula for Dn. d) Hence, determine how many months are need to repay the debt if "x" = 12,000. 2. A partion of "n" is a sequence of positive integers (λ1,λ1..,λr) such that λ1 ≥ λ2 ≥ .. ≥ λr and λ1 + λ2 + .. + λr = n. For example, let "n" = 5, then (5), (4,1), (3,2), (3,1,1), (2,2,1), (2,1,1,1), (1,1,1,1,1) are all the partions of 5. So we say the number of partions of 5 is 7. Derive a formula for the number of partions of a general "n". 3. A binary tree is a tree where each node can have up to two child nodes. The figure below is an example of a binary tree. a) Let cn be the number of unique arrangements of a binary tree with totally n nodes. Let C(z) be a generating function of cn. b) Let formula_73 be a power series. Hint: Instead of doing recursion of finding the change in cn when adding nodes at the buttom, try to think in the opposite way, and direction.(And no, not deleting nodes) Project - Exponential generating function. This project assumes knowledge of differentiation. (Optional)0. 1. Consider formula_75 (Optional)2. 3. The function E(x) is the most fundamental and important exponential generating function, it is similar to the ordinary generating function, but with some difference, most obviously having a fractorial fraction attached to each term. 4. Apart from A(x) defined in question 2, let formula_78 Notes: Question with *** are difficult questions, although you're not expected to be able to answer those, feel free to try your best. 

Grammatica 4-1 ~ The indefinite articles "een" en "geen". In the previous lesson you were introduced to the definite articles—'the' in English and het or de in Dutch. Indefinite articles precede nouns in the same way that definite articles do, but convey a general or indefinite sense. These are 'a' or 'an' in English. Thus, 'the book' or "het boek" refers to a definite or specific book, whereas 'a book' or "een boek" is indefinite about which book is referred to. If your mother tongue does not have articles (such as Russian, Polish etc.) try this. Dutch indefinite articles only come in one form ("een"), so they don't display gender. The use of definite and indefinite articles is virtually the same as in English. The few deviations are best learned when listening to the language or speaking it. Please note (see also previous lesson) that the indefinite article has the same form as the numeral one ("één"). One could argue that one is a clitic form of the other. To denote the difference, one could place accents on the numeral. Also, there is a difference in pronunciation. The numeral "één" (one) is pronounced /e:n/, the article "een" (a) with a much weaker /ən/. Occasionally Dutch has one and English the other: Notice that "one" is used here in the meaning of "a certain", not say in contrast to two or three. There is an inflected form "ene" that is used independently: Negation. In English a negative of an indefinite article is simply formed by adding not: Alternatively one can drop the article and say: In Dutch there is a special negative of een: "Geen" is used both with singular and with plural indefinite nouns Relation to "niet". The combination niet + een is only used in contrasting things: With definite nouns negation involves adding "niet", usually at end: Notice also that Dutch does not use the auxiliary verb "to do" for negations: Relation to "wel". When a negative is negated Dutch uses the adverb "wel" to express that. English has to use a construction involving the verb "to do": This adverb is also used for contrasting: Notice that the second part of the Dutch sentence does not even have a verb. Dutch is quite an 'adverbial' language. If the adverb expresses the meaning sufficiently, why bother with superfluous things like verbs? A game of "Yes, you do" - "No, I don't" sounds like "Welles!" - "Nietes!" in Dutch Much like in French, there are two words for "yes". A simple confirmation is "ja" (French: "oui"); a negation of a negative is "jawel" (French: "si") Grammatica 4-2 ~ More pronouns. Possessives. Recall the following from Gesprek 3-1: Which translates as: The sentence demonstrates one of the possessive pronouns. In the singular these are The neuter "zijn" for "its" is not used very much in Dutch, as we shall see in lesson 8 it often gets replaced by "ervan". The above pronouns like "jouw" often turn into a weak (clitic) form je that is used when the emphasis is on something else, such as the motorcycle in this case. In the spoken language this holds for all the ones shown in the table, but in the written language "je" is the most generally accepted clitic. When written they are: Dutch does not have a possessive case as English does. In English one could say "this house of mine", where mine (and yours, hers, his, ours, yours, theirs) is possessive case. Dutch uses objective case for this: "dit huis van mij" as if 'van' ("of") is a preposition. However, for a sentence like "is this yours or his?" Dutch would use nominalized pronouns (pronouns turned into nouns) with an inflection -e usually accompanied by a definite pronoun "de" or "het": The possessive for the polite 2nd person pronoun "u" is "uw" and in the south this same word is used to refer to the "gij" pronoun (but no formal usage is implied). "Uw" does not have a clitic form and the same can be said about the possessives of the plural. Uw does have an inflected form "uwe" used for nominalization. For the third person plural the possessive is "hun" and it follows the same pattern: For the first person plural Dutch has "ons". This is the only one that follows the rules for inflection of the adjectives. I.e. its inflected form is used for de-words but not for het-words: The inflected form is once again also used for nominalization: For the second person plural the possessive is "jullie" which cannot be inflected and cannot be used as a noun, which necessitates a construction involving "van": See Dutch/Appendix 3 for a table of the possessive pronouns. Quizlet. Practice the possessive pronouns at Quizlet (24 term, includes demonstrative and reciprocal pronouns) Demonstratives. In English, this is used as demonstrative pronoun to indicate something in proximity. "That" indicates greater distance. In Dutch a similar distinction exists, but gender plays a role: So, one replaces 'de' by deze and 'het' by dit. At a greater distance: Notice that often when English has th, Dutch will have d: A third, even more distant pronoun exists (gene, gindse), but it is about as common as its English equivalent "yon", "yonder". Again, the two languages betray their kinship. In some words, a g in Dutch corresponds to a y in English.. Compare: Using demonstrative pronouns instead of personal pronouns. Recall: As we have seen Dutch is on its way to a two-gender system. For inanimate nouns, this makes demonstrative pronouns a more attractive choice to refer things by than personal pronouns. Compare: As you see demonstratives do not distinguish whether a word is feminine or masculine and follow the same common-neuter pattern as the articles. Compare: Note: because "de auto" is not neuter, it is not correct to say: "Het" is duur. But saying "hij" is duur or "zij" is duur makes the word specifically masculine or feminine. Using die avoids the issue, because die follows the "common gender" pattern of the definite article. Increasingly, personal pronouns are reserved for reference to persons (natural gender as in English). To refer to things people resort to substituting the demonstratives. Reflexive and Reciprocal pronouns. In English reflexive pronouns always carry the ending -self -selves: myself, themselves etc. In Dutch that is not always so. In fact, for a verb that is "always" reflexive, like "zich vergissen" (to be mistaken) the ending "cannot be used": In other words the reflexive pronoun is identical to the clitic object form of the personal pronoun, except in the third person where it is zich. The pronoun u was originally a third person (It stems from U.E. uwe edele, something like: your nobility, your honor) which explains the zich for this pronoun. Vergissen can only be used with zich, but some verbs can be used with or without a reflexive pronoun. In that case -zelf may be added: Ik was me/mij/mezelf/mijzelf. - I wash myself. This topic is revisited in Lesson 16 The most important reciprocal pronoun is elkaar - each other Grammatica 4-3 Plural of nouns. We already have seen some things about the plural above: Forming the plural of the noun itself is a bit more complicated. In English it is basically always done with -s, but in Dutch that is different Recall: "...tafels en stoelen..." With few exceptions like "ox - oxen" pretty much all words simply get an -s in English. Dutch however has two main ways to form a plural: by adding -s and by adding -en. The latter is pronounced /-ən/, /-ə/ or even as a syllabic /-n/ depending on the region. Which plural applies is best learned case by case as gender is, although we can attempt a general rule: The ones in -a, -o, -i and -y get an apostrophe before the -s Unfortunately there are lots of exceptions. Many recent (latinate) loans from English or French and all diminutives get a -s. Words in -te and -aar usually get -s: Amongst the many words that get -en are the ones in -ing: Vowel changes. Most monosyllabic words have -en in the plural: In the latter case, notice that one of the a's is dropped in the spelling of the plural. This difficulty is related to the fact that most Dutch vowels occur in two varieties, a closed one and an open one. Dutch spelling has a rather ingenious and systematic way of denoting which one is intended. It involves the doubling of either vowels or consonants. Compare: In this case the vowels remain the same in the plural, but notice the doubling: It is customary to call the first sound [ɔ] a 'short o' and the second [o] a 'long o', but this terminology can be rather confusing. There are languages like Czech or Gàidhlig where vowels are indeed distinguished purely on their length. In Dutch, however, the difference in length ("quantity") is actually pretty negligible, but the difference in vowel sound ("quality") is not. This presents a problem for speakers of the many languages with a five-vowel system, like Italian, Russian, Arabic or isiXhosa whose ears are not accustomed to this kind of difference. Anglophones usually do quite well. The following five vowels possess open ('long') and covered ('short') varieties: Vowel changes. For non-native speakers a complication arises in those cases where the actual vowel changes ('lengthens') in the plural, compare: The vowel /ɑ/ in pad and padden is approximately as in father. Paden has a vowel /a/ like in broad American 'Oh, my God' (In Dutch the spelling would be: Gaad). Also, notice the gender difference of the two words. Vowel change is systematic in the plural of the past of certain strong verbs (class 4 and 5; see 6). A few words show vowel changes other than between the open and closed variety of the same vowel: Words ending in -heid get -heden: -eren. There are about a dozen plurals in Dutch that end in -eren: The ending -eren is essentially a double plural. It derives from a plural in -er and in some compounds that is still visible: -ën. Some words in -ie have an -en plural that requires a diaeresis ("trema" in Dutch). The spelling depends on where the stress falls: Notice that in Dutch orthography the stress of a word cán be indicated with an acute accent, but this is only permitted if otherwise ambiguity might arise. A "trema" (diaeresis) is also used after -ee: Changes of consonants. If the root of the word ends in a -z of -v this is both written and pronounced as -s and -f in the singular, but the voiced consonant returns in a plural on -en: Such words typically have a 'long' vowel ("aa" in this case) or diphthong. With "short" vowel the consonant typically remains voiceless and is doubled in the plural Lieden, lui. Some compounds of -man have a plural involving "lieden" or "lui" (people) Classical plurals. Occasionally a Latin or Greek plural is preserved in Dutch: Woordenlijst 4. &lt;br clear="all"&gt; Quizlet. The vocabulary or this lesson can be practiced at Quizlet (28 terms) Progress made. Cumulative count &lt;br clear="all"&gt; Quizlet. The vocabulary or this lesson can be practiced at Quizlet (28 terms) Progress made. If you have thoroughly studied the above lesson you should Cumulative vocabulary count &lt;br clear="all"&gt; 

 Lektion 6 "Die Wohnung" ~ The Apartment Gespräch 6-1 ~ Ein Bruder besucht Markus. This incomplete story and conversation introduces terms for items around the house (or apartment). Vokabeln 6-1.  der Bruder brother  die Eltern parents  die Küche kitchen  das Schlafzimmer bedroom  die Vorlesung class, instruction (at a university)  die Wohnung apartment, flat  das Wohnzimmer living room  das Zimmer, die Zimmer room(s)  es gibt there is  gegen Abend towards evening  gern haben like (i.e., "to gladly have")  Herein! Come in!  sich umsehen look around  zeigen show  besuchen visit, attend (classes)  grüßen greet  mieten rent  sein his (a possessive adjective) Grammatik 6.1 ~ Introduction to verb conjugations. In German, every grammatical person has, or potentially has, its own unique verb form. Describing the various verb forms is called verb conjugation. This variation in verb form is certainly one of the things that makes German grammar somewhat difficult for English speakers to learn. In English, only the 3rd person singular might differ from the verb form used with all of the other persons (see Grammatik 1-3) and that difference is made by adding an ending of 's' or 'es'. For example: I/you/we/they 'go', but he/she/it 'goes'. Let us have a closer look at German verbs. Usually, the infinitive form of a verb in German ends with -en—for examples, consider these verbs you have already learned: "gehen ('go'), "haben ('have'), and "studieren" ('study'). In order to "build" the different verb forms (that is, conjugate a verb), first cut off the '-en' ending from the infinitive. Then append a new ending according to the grammatical person. For regular verbs it works essentially as follows: As you see in this example using the verb "gehen", the singular 1st person ends with -e, the 2nd person with -st and 3rd person (no matter what gender) ends with -t. As for the plural forms, note that 1st and 3rd person in plural number (see Grammatik 1-3) are built the same way as the infinitive. Again note that, in English, only the verb form for the 3rd person singular is "unique". An easy way to remember the regular verb endings is the following mnemonic "Elephants standing together enjoy trumpeting endlessly". Seems simple enough. However, realize we are discussing here only the regular verb forms in the present tense ("Präsens"). You will learn quite soon that, unfortunately, there are many exceptions from these simple rules. An important one is the irregular verb "sein" ('to be') which is irregular in English as well (I am, you are, he is...). At least 1st and 3rd person plural are the same. Another important verb is "haben" ('to have'): You see, it's not too irregular—only the 2nd and 3rd person singular constitute a small exception since the 'b' has vanished. English is somewhat curious in this respect as well: 'I have', but 'he has'. Future lessons will introduce you to the many irregular verbs in German. But you should now recognize what is happening to the verbs in German sentences. They are reflecting the person and number of their nominative case subjects. Recall these sentences from past lessons (verbs underlined here):  "Danke, es geht mir gut" Thanks, it goes well with me (verb is "gehen")  "Ich habe viel Arbeit" I have much work (verb is "haben")  "Ist er zu Besuch?" Is he visiting? (verb is "sein")  "Du bist ein Schwein!" You are a pig! (verb is "sein")  "Wie heißen Sie?" What are you called? (verb is "heißen", and pronoun is formal)  "Wir spielen eine Stunde lang!" We play for one hour! (verb is "spielen")  "Sie liegt am Ausfluss des Zürichsees." It lies at the outlet of Lake Zurich (verb is "liegen") Grammatik 6.2 ~ Case in German nouns. Through our discussions on the personal pronouns, you have learned how pronouns have case. Nouns also have case—and in German, noun case can be expressed by the definite article ("der"). Recall this table from Lektion 3: These "der"-words reflect noun gender in the nominative case—appropriate whenever a noun is used as the subject of a sentence. For other cases, the "der" words change. Expanding the table to present nominative (NOM.), accusative (ACC.), dative (DAT.), and genitive (GEN.) cases: Note, there are also "der"-word forms to be used for plural nouns. Fortunately, these are the same, no matter what the gender of the singular noun. For future reference, you can find the "der"-words summarized in Anhänge Drei. The following examples demonstrate the use of the definitive article in various parts of speech:  "Du hast die Wurst und den Käse." You have the sausage and the cheese.  (accusative case)  "Die Geschäftsleute verstehen die Arbeit" The business associates understand the work.  (nominative and accusative cases)  "Sie liegt am Ausfluss des Zürichsees." It lies at the outlet of (the) Lake Zurich.  (genitive case)  "Zürich ist die größte Stadt der Schweiz." Zurich is the largest city in (of the) Switzerland.  (nominative and genitive cases) In the last example, remember that in both English and German, the noun (or pronoun) that follows the verb 'to be' is a predicate noun, for which the correct case is the nominative. That is why, in English, 'It is I' is grammatically correct and 'It is me' is incorrect. So consider the following (and note that case of each definite article is the same as in the last example above):  "Zürich ist der Kanton der gleichnamigen Stadt." Zurich is the canton of the same named city. Grammatik 6.3 ~ Commands.  "Ruf sie an, bitte!" Call her, please.  or "Ruf sie bitte an!"  "Gehen Sie nach Hause!" Go home (formal).  "Kommt mit!" Come with (plural)!  "Gib es mir!" Give me it! Notice that in these sentences there are no subjects (except for #2). In German, as in English, there is a "commandative form", a way to demand something using an understood you. In English, there is only one you-form and one command form. In German, since there are three you's, there are three ways to command. If the subject is singular ("du"), then the verb has no ending. If it is irregular, it takes the du-form, such as in essen (Iss!) or lesen (Lies!). If there is a plural subject ("ihr"), then the verb takes the ihr-form. Nothing else is changed. Most of the time, ihr-commands are used with children, but that is not always the case. In both of these sentences, the du or ihr is omitted. Formal is normal. The "Sie" stays (after the verb) and the verb is in its formal form. Although it is worded like a question, in written or spoken form, it is easy to tell the difference. 

Moles. The mole is a unit used to measure the number of particles in a substance (where particles can be atoms, molecules, electrons, ions etc.). One mole of substance contains 6.02205 × 1023 particles. This rather large quantity is a constant called "Avogadro's Number", which is abbreviated as "NA". It simply refers to the number of particles in one mole. formula_1 For example, one mole of helium gas contains 6.022×1023 atoms; 2 moles of oxygen gas contains 1.204×1024 molecules. Moles and Mass. It would be useful to be able to convert between the number of moles and the mass of a substance. The "molar mass" is the mass of one mole of a substance, measured in grams per mole or formula_2. Not by coincidence, the "molar mass" is the same as the "atomic weight" of a substance (this is how the concept of moles was originally defined). For example, the molar mass (and atomic weight) of carbon-12 is 12 g·mol−1, so 1 mole of carbon-12 weighs 12 grams. Similarly, chlorine has a molar mass of 35.45 g·mol−1, so one mole of chlorine atoms weighs 35.45 g. There are two chlorine atoms in one molecule of Cl2 gas, so one mole of Cl2 gas weighs 70.90 g. Measurement and Quantities. "See the Units Appendix for tables of SI fundamental and derived units. 

Detailed Objective. (LPIC-1 Version 5.0 Weight: 3 Description: &lt;br&gt; Candidates should be able to perform package management using the Debian package tools. Key knowledge areas: The following is a partial list of the used files, terms and utilities: Package Structure. In order to understand how to use Debian's package management system, it would be useful to first have an understanding of how a Debian package is named. For example, the package ncftp_3.1.3-1_i386.deb has 5 major parts: Note that there is special significance to the use of underscores(_) and hyphens(-); an underscore shall separate the name of the program and its version, a hyphen shall separate a version number and the revision number, and an underscore shall separate the revision number and the architecture. dpkg. dpkg is the "grandaddy" or "back-end" of the Debian Package Management System. Features present in the more advanced tools are not available to dpkg but it is nevertheless a useful tool. Some notes: Some of the more common functions used by administrators by dpkg are:&lt;br&gt; Adding, Removing, and Configuring Packages Querying Package Information Updating Package Information dpkg-reconfigure. dpkg-reconfigure reconfigures packages after they have already been installed. Dselect. The utility that allows you on debian to easely add/remove packages is dselect. There is on dselect an interactive menu that will allow you to install/remove packages. Care must be taken with this utility. You can damage your system. Dselect menu example:  Debian `dselect' package handling frontend.  0. [A]ccess Choose the access method to use.  1. [U]pdate Update list of available packages, if possible.  2. [S]elect Request which packages you want on your system.  3. [I]nstall Install and upgrade wanted packages.  4. [C]onfig Configure any packages that are unconfigured.  5. [R]emove Remove unwanted software.  6. [Q]uit Quit dselect.  $ dselect - list of access methods  Abbrev. Descriptio cdrom Install from a CD-ROM.  * multi_cd Install from a CD-ROM set.  nfs Install from an NFS server (not yet mounted).  multi_nfs Install from an NFS server (using the CD-ROM set) (not yet mounted).  harddisk Install from a hard disk partition (not yet mounted).  mounted Install from a filesystem which is already mounted.  multi_mount Install from a mounted partition with changing contents.  floppy Install from a pile of floppy disks.  apt APT Acquisition [file,http,ftp] apt-get. If you know the name of a package you want to install, use apt-get. You must configure the sources.list file. This same file is used when you choose the apt access method of dselect. The location is /etc/apt. Apt-cache. To find the name of a package that you want to install use apt-cache. apt-cache main options are :  user@host:~$ apt-cache search gimp  babygimp - An icon editor in Perl-Tk  blackbook - GTK+ Address Book Applet  cupsys-driver-gimpprint - Gimp-Print printer drivers for CUPS  escputil - A maintenance utility for Epson Stylus printers  filmgimp - A motion picture editing and retouching tool Resources. APT HOWTO&lt;br&gt; http://www.debian.org/doc/manuals/apt-howto/index.en.html&lt;br&gt; dselect Documentation for Beginners&lt;br&gt; http://www.debian.org/doc/manuals/dselect-beginner/&lt;br&gt; 

For any triangle with vertices formula_1 corresponding angles formula_1 and corresponding opposite side lengths formula_3 , the Law of Sines states that Each of these expressions is also equal to the diameter of the triangle's circumcircle (the circle that passes through the points formula_1). The law can also be written in terms of the reciprocals: Proof. Dropping a perpendicular formula_7 from vertex formula_8 to intersect formula_9 (or formula_9 extended) at formula_11 splits this triangle into two right-angled triangles formula_12 and formula_13 . We can calculate the length formula_14 of the altitude formula_7 in two different ways: By using the other two perpendiculars the full law of sines can be proved. QED. Application. This formula can be used to find the other two sides of a triangle when one side and the three angles are known. (If two angles are known, the third is easily found since the sum of the angles is formula_21 .) See Solving Triangles Given ASA. It can also be used to find an angle when two sides and the angle opposite one side are known. Area of a triangle. The area of a triangle may be found in various ways. If all three sides are known, use Heron's theorem. If two sides and the included angle are known, consider the second diagram above. Let the sides formula_22 and formula_23 , and the angle between them formula_24 be known. The terms /alpha and /gamma are variables represented by Greek alphabet letters, and these are commonly used interchangeably in trigonometry just like English variables x, y, z, a, b, c, etc. From triangle formula_25 , the altitude formula_26 is formula_27 so the area is formula_28 . If two angles and the included side are known, again consider the second diagram above. Let the side formula_23 and the angles formula_24 and formula_31 be known. Let formula_32 . Then Thus 

Law of Cosines. The Pythagorean theorem is a special case of the more general theorem relating the lengths of sides in any triangle, the law of cosines: where formula_2 is the angle between sides formula_3 and formula_4 . Does the formula make sense? This formula had better agree with the Pythagorean Theorem when formula_5 . So try it... When formula_5 , formula_7 The formula_8 and the formula reduces to the usual Pythagorean theorem. Permutations. For any triangle with angles formula_9 and corresponding opposite side lengths formula_10 , the Law of Cosines states that Proof. Dropping a perpendicular formula_14 from vertex formula_15 to intersect formula_16 (or formula_16 extended) at formula_18 splits this triangle into two right-angled triangles formula_19 and formula_20 , with altitude formula_21 from side formula_22 . First we will find the lengths of the other two sides of triangle formula_19 in terms of known quantities, using triangle formula_20 . Side formula_22 is split into two segments, with total length formula_22 . Now we can use the Pythagorean Theorem to find formula_4 , since formula_33 . The corresponding expressions for formula_3 and formula_22 can be proved similarly. The formula can be rearranged: and similarly for formula_37 and formula_38 . Applications. This formula can be used to find the third side of a triangle if the other two sides and the angle between them are known. The rearranged formula can be used to find the angles of a triangle if all three sides are known. See Solving Triangles Given SAS. 

Database Programming. Introduction. A database is a repository of information managed by a database engine which ensures integrity of data and fast access to the data. A very common implementation of a database is a Relational Database Management System (RDBM). To users, the information in a database can be accessed by using Structured Query Language (SQL) a database language common to most databases. However, SQL only provides commands to access and manipulate the data in the database. For any complex application, there commonly is a need for conditions, branching, and loops. None of these are a part of the SQL language. In order to fill this gap, many common programming languages allow integration with SQL through a common library such as Open Data Base Connectivity (ODBC), Object Linking and Embedding (OLE), and sometimes with application programming interfaces or libraries supplied with the database. In addition, most databases now have a simple language of their own which allows simple control for applications which do not need the full power of standard languages like C++ and Pascal. These simple languages are used to write stored procedures and are proprietary to each database. An alternate approach taken by other languages like is to have a data model that includes persistent variables that are shared among multiple processes. This approach hides the database operations within the programming language instead of hiding the language within the database. 

Date: Fri, 16 Jan 2004 14:33:06 -0500&lt;br&gt; From: "Wilma Silbermann" &lt;wilma@lpi.org&gt;&lt;br&gt; To: "m strawser" &lt;mind3ras3r@yahoo.com&gt;&lt;br&gt; Subject: Re: reprinting of exam objectives&lt;br&gt; Hi there Mark, We are pleased that you are creating more training material for the LPI exams and would have no problem approving your request to reprint the objectives as long as you ensure you show that LPI does not approve or endorse your book in any way. You may use the logo in the same manner - please read the logo license; https://www.lpi.org/en/logo.html Please give me a call if you have any further concerns or questions on this or any other matter. Thanks again for your support of LPI. Wilma 

Swing is a tool kit in Java which provides a way to build cross platform user interfaces. It is built on top of and designed as a replacement for AWT, the other UI toolkit built into Java. SWT is a third toolkit you may hear about. SWT is an open source toolkit and is a full topic in of itself, see SWT's Homepage for more information. Overview. Swing provides many controls and widgets to build user interfaces with. Swing class names typically begin with a J such as codice_1, codice_2, codice_3. This is mainly to differentiate them from their AWT counterparts and in general are one-to-one replacements. Swing is built on the concept of "Lightweight components" vs AWT and SWT's concept of "Heavyweight components". The difference between the two is that the Lightweight components are rendered (drawn) using purely Java code, such as drawLine and drawImage, whereas Heavyweight components use the native operating system to render the components. Some components in Swing are actually heavyweight components. The top-level classes and any derived from them are heavyweight as they extend the AWT versions. This is needed because at the root of the UI, the parent windows need to be provided by the OS. These top-level classes include codice_4, codice_3, codice_6 and codice_7. All Swing components to be rendered to the screen must be able to trace their way to a root window or one of those classes. So what does using Swing get you? So far we've only talked about components and rendering. Well, to start with you get the following. 

Parliament. Parliament is the supreme law-making body in the United Kingdom. It is made up of two "Houses of Parliament", namely the "House of Commons" and the "House of Lords", as well as the Sovereign. The Sovereign's involvement in the life and working of Parliament is purely formal. In constitutional theory, Parliament in its strictest sense is sometimes referred to as the "Queen-in-Parliament"; this contrasts with the more ordinary use of the term "Parliament", meaning just the two Houses of Parliament. Within the British constitutional framework, the Queen-in-Parliament is supreme ("sovereign"), able to make, alter, or repeal any law at will. Both Houses of Parliament meet at the "Palace of Westminster". Parliaments and Sessions. As with most legislatures, Parliament does not continue in perpetual existence. Typically, the "life" of a Parliament is around four years. Parliament is initially "summoned" by the Sovereign. This now always occurs after there has been a general election. Once assembled, and a Speaker has been chosen by the House of Commons, Parliament is formally "opened" by the Sovereign. The business of the two Houses is arranged into "sessions", which usually last a year (running from around October or November each calendar year). However, there is usually a long "recess" during the summer months, when business is temporarily suspended. The opening of each parliamentary session is conducted in accordance with a great deal of traditional ceremony. The Sovereign takes his or her seat on the throne situated in the chamber of the House of Lords, and the Gentleman Usher of the Black Rod (one of that House's officers) is commanded to summon the House of Commons. When Black Rod reaches the door of the Commons, it is slammed shut in his face, to symbolise the right of the Commons to debate without royal interference. Black Rod then solemnly knocks on the door with his staff of office; on the third knock, the door is opened, and he is permitted to enter and deliver his message. MPs then proceed from the Commons to the House of Lords, to hear the "Speech from the Throne", more commonly known as the "Queen's Speech". The Speech outlines the Government's legislative proposals for the session; while worded as if it's the Sovereign's own policy, the Speech is in fact entirely drafted by Government ministers. Each session is ended by a "prorogation". The Commons are formally summoned to the House of Lords, where another formal Speech is read out, summing up the work of the two Houses of Parliament over the course of the session. In practice the Sovereign no longer attends for the prorogation; "Lords Commissioners" are appointed to perform the task, and one of their number also reads out the Speech. By law, each Parliament must come to an end no later than five years from its commencement; this is known as "dissolution". Because a dissolution is necessary in order to trigger a general election, the Prime Minister was effectively able to choose to hold elections at a time that seems the most advantageous to his or her political party. After the 2010 General Election the Liberal Democrat party made the introduction of a law on fixed term parliaments a requirement for forming a coalition government with the Conservative party. A coalition was necessary as the result of the election meant no party had an overall majority of seats. The main opposition party from the previous parliament, the Conservative party, gained the most seats and under parliamentary protocol had the first opportunity to try and form a government. They concluded a formal deal with the Liberal Democrats to govern together. The Liberal Democrats had insisted on fixing the term of parliaments to reduce the inherent advantage the governing party had in being able to choose the moment to hold an election. There are provisions in the Fixed Term Parliaments Act to allow an early election with the consent of Parliament. These provisions were used in 2017. Although the duration of Parliament has been restricted to five years since 1911, legislation was passed during both World Wars to extend the life of the existing Parliament; this meant that the Parliament summoned in 1935 eventually continued in existence for around ten years, until 1945. House of Commons. Composition. While sometimes described as the "lower house", the House of Commons is by far the most important of the two Houses of Parliament. Members of the House of Commons are known as "Members of Parliament", or "MPs". The entire United Kingdom is subdivided into "constituencies", each of which returns one MP to sit in the House of Commons. There are presently 650 constituencies, however the exact number fluctuates over time as the boundaries of constituencies are periodically reviewed by Boundaries Commissions set up for each part of the UK. Constituencies are intended to have roughly equal numbers of voters, but in practice the smallest and largest constituencies can have a significant difference in size. At each "general election" all seats in the House of Commons become vacant. If a seat becomes vacant during the life of a Parliament (i.e. between general elections), then a "by-election" is held for that constituency. The election for each constituency is by secret ballot conducted according to the "First-Past-the-Post" system: the candidate with the most votes is "returned" as MP. Qualifications of voters. A person must be aged at least eighteen in order to vote. The following nationalities are entitled to vote at parliamentary elections: Irish and Commonwealth citizens must have been resident in the United Kingdom. British citizens who are resident abroad are only able to vote if they had been resident in the United Kingdom within the previous 15 years. Certain categories of people are unable to vote: By convention, close relatives of the Sovereign also do not vote. Qualifications of MPs. Anyone who is not disqualified to vote is also qualified to be an MP, except the following: Resignation as an MP. Since the 17th century, the House of Commons has asserted that MPs may not resign. However, in practice members are able to resign by the legal fiction of appointment as "Crown Steward and Bailiff of the three Chiltern Hundreds of Stoke, Desborough, and Burnham", or as "Crown Steward and Bailiff of the Manor of Northstead". Neither of these offices carries any duties, but have been preserved in force so that those appointed to them automatically lose their seats in the House of Commons as having accepted an "office of profit under the Crown". Speakership and procedure. The House of Commons is presided over by the "Speaker". There are also three Deputy Speakers, with the titles of "Chairman of Ways and Means", "First Deputy Chairman of Ways and Means", and "Second Chairman of Ways and Means". The Speaker and his or her deputies are elected at the commencement of a Parliament, and serve until its dissolution. Following a general election, the "Father of the House" (the member with the longest unbroken service in the House, who is not also a Minister of the Crown) takes the chair. If the Speaker from the previous Parliament has been returned as a member of the new Parliament, and intends to continue in office, then the House votes on a motion that the member take the chair as Speaker. Otherwise, or if the motion for his or her re-election fails, then members vote by secret ballot in several rounds; after each round, the candidate with the fewest votes is eliminated. The election ends when one member secures a majority of votes in a particular round. Thereafter, the Speaker-elect leads the House of Commons to the House of Lords, where the Lords Commissioners (five Lords representing the Sovereign) officially declare the Royal Approbation (approval) of the Speaker, who immediately takes office. The Speaker traditionally lays claim to all of the House's privileges, including freedom of speech in debate, which the Lords Commissioners then confirm on behalf of the Sovereign. If a Speaker should choose to resign from his post during the course of Parliament, then he must preside over the election of his successor. The new election is otherwise conducted in the same manner as at the beginning of a Parliament. The new Speaker-elect receives the Royal Approbation from Lords Commissioners; however, the ceremonial assertion of the rights of the Commons is not repeated. The Speaker is expected to act impartially. He or she is an important figure within the House of Commons, controlling the flow of debate by selecting which members get to speak in debates, and by ensuring that the customs and procedures of the House are complied with. The Speaker and his deputies do not generally speak during debates, nor vote at divisions. The Speaker also exercises disciplinary powers. He or she may order any member to resume his or her seat if they consistently contribute irrelevant or repetitive remarks during a debate. An individual who has disregarded the Speaker's call to sit down may be requested to leave the House; if the request is declined, then the Speaker may "name" the member. The House then votes on whether to suspend the member in question for a certain number of days, or even, in the case of repeated breaches, for the remainder of the session. In the most serious cases, the House may vote to expel a member. In the case of grave disorder, the Speaker may adjourn the House without a vote. The House votes on all questions by voice first. The Speaker asks all those in favour of the proposition to say "Aye," and those opposed to say "No". The Speaker then assesses the result, saying "I think the Ayes have it" or "I think the Noes have it", as appropriate. Only if a member challenges the Speaker's opinion is a "division", or formal count, called. During a division, members file into two separate lobbies on either side of the Commons chamber. As they exit each lobby, clerks and tellers count the votes and record the names. The result is then announced by the Speaker. In the event of a tied vote, the Speaker (or other occupant of the Chair) has a "casting vote"; however, conventions exist that the Speaker would cast a vote to maintain the "status quo." In effect moving bills on to further scrutiny, but not pass a bill into law. House of Lords. Composition. Generally speaking, membership of the House of Lords is by appointment for life. However, up until 1999, hereditary peers were also members of the Lords; when this right was abolished, a compromise measure allowed them to elect ninety of their number to continue as members. Certain office-holders are also ex officio members of the House of Lords: The Earl Marshal and Lord Great Chamberlain are mostly ceremonial offices. In addition to the three ex officio bishops, the 21 longest-serving diocesan bishops also sit in the Lords. The general qualifications for sitting and voting in the Lords are: Speakership and procedure. The "Lord Speaker" is elected by the House. Until recently his or her duties were carried out by the Lord Chancellor, a Minister of the Crown. In contrast with the Speaker of the House of Commons, the Lord Speaker has a relatively minor role, since the House of Lords is generally self-governing: the House itself decides upon points of order and other such matters. The seat used by the Lord Speaker is known as "the Woolsack". Similar to the House of Commons, the Lords also vote by voice first. The Lord Speaker (or whoever else is presiding) puts the question, with those in favour saying "Content," and those opposed saying "Not-Content." If the Lord Speaker's assessment of the result is challenged, a division follows, with members voting in the appropriate lobby just as is done in the Commons. The officer presiding may vote from his or her place in the chamber rather than from a lobby. In the case of a tie, the result depends on what type of motion is before the House. A motion that a bill be advanced to the next stage or passed is always decided in the positive, while amendments to bills or other motions are decided in the negative, if there is an equality of votes. Acts of Parliament. Legislation passed by Parliament is in the form of an "Act of Parliament". A draft law is known as a "Bill". A bill passes into law provided that it has either been passed by both Houses of Parliament, or the provisions of the Parliament Acts have been complied with; and provided it has received the Royal Assent. A bill must pass through several stages in both of the two Houses. A bill is "read" three times in each House. The First Reading for Public Bills is almost always a formality. The Second Reading is a debate on the merits of the general principles behind the bill. Next follow the Committee and Report stages. The Third Reading is a vote upon the bill as a whole, as amended during the Committee and Report stages. Once the House into which the bill was first introduced has finished with it, the bill is then introduced into the other House. Any amendments by the second House then have to be agreed to by the first before the bill can proceed. Bills are classified as either "Government Bills" or as "Private Members' Bills". Ministers of the Crown introduce Government Bills; private members introduce Private Members' Bills. Bills are also classified as "Public", "Private", "Personal" or "Hybrid". Public bills create laws applied generally (for instance, reforming the nation's electoral system). Private bills affect a specific named company, person or other entity (for instance, authorising major constructions on specific named public lands). Personal bills are private bills that confer specific rights to specific named individuals (for example by granting the right to marry a person one would not normally be allowed to wed). Hybrid bills are public bills that directly and specially affect private interests. Public Bills. A Public Bill's First Reading is usually a mere formality, allowing its title to be entered in the Journals and for its text to be printed by the House's authority. After two weeks, one of the bill's supporters moves "that the bill be now read a second time". At the second reading debate, the bill's general characteristics and underlying principles, rather than the particulars, are discussed. If the vote on the Second Reading fails, the bill dies. It is, however, very rare for a Government bill to be defeated at the Second Reading; such a defeat signifies a major loss. In the House of Commons, following the Second Reading, various procedural resolutions may need to be passed. If the bill seeks to levy or increase a tax or charge, then a "Ways and Means Resolution" has to be passed. If it involves significant expenditure of public funds, then a "Money Resolution" is necessary. Finally, the government may proceed with a "Programme Motion" or an "Allocation of Time Motion". A Programme Motion outlines a timetable for further debate on the bill and is normally passed without debate. An Allocation of Time Motion, commonly called the "Guillotine", limits time available for debate. Normally, a Programme motion is agreed to by both parties while an Allocation of Time Motion becomes necessary if the Opposition does not wish to cooperate with the Government. In the House of Lords, there are no Guillotines or other motions that limit the time available for debate. Next, the bill can be committed to a committee. In the House of Commons, the bill may be sent to the "Committee of the Whole House", a "Standing Committee", a "Special Standing Committee" or a "Select Committee." The Committee of the Whole House is a committee that includes all members of the House and meets in the regular chamber. The Speaker is normally not present during the meetings; a Deputy Speaker normally takes the chair. The procedure is used for parts of the annual Finance Bill and for bills of major constitutional importance. More often, the bill is committed to a Standing Committee. Though the name may suggest otherwise, the membership of Standing Committees is temporary. There can be from sixteen to fifty members; the strength of parties in the committee is proportional to their strengths in the whole House. It is possible for a bill to go to a Special Standing Committee, which is like a Standing Committee except that it may take evidence and conduct hearings; the procedure has not been used in several years. Finally, the bill may be sent to a Select Committee. Select Committees are permanent bodies charged with the oversight of a particular Government department. This last procedure is rarely used; the quinquennial Armed Forces Bill, however, is always referred to the Defence Select Committee. In the House of Lords, the Bill is committed to the "Committee of the Whole House", a "Public Bill Committee", a "Special Public Bill Committee", a "Select Committee" or a "Grand Committee". The most common committee used is the Committee of the Whole House. Sometimes, the bill is sent to a Public Bill Committee of twelve to sixteen members (plus the Chairman of Committees) or to a Special Public Bill Committee of nine or ten members. These committees correspond in function to the Commons Standing and Special Standing Committees, but are less often utilised. Select Committees may also be used, like in the Commons, though it is rare for this to be done. The Grand Committee procedure is the only one unique to the House of Lords. The procedure is reserved for non-controversial bills that must be passed quickly; a proposal to amend the bill is defeated if a single member votes against it. In both Houses, the committee used considers the bill clause-by-clause and may make amendments. Thereafter, the bill proceeds to the "Consideration" or "Report Stage". This stage occurs on the Floor of the House and offers it an opportunity to further amend the bill. While the committee is bound to consider every single clause of the bill, the House need only debate those clauses which members seek to amend. Following the Report Stage, the motion "that the bill be now read a third time" is considered. In the House of Commons, there is a short debate followed by a vote; no further amendments are permitted. If the motion passes, then the Bill is considered passed. In the Lords, however, amendments may be moved. Following the vote on the third reading, there must be a separate vote on passage. After one House has passed a bill, it is sent to the other for its consideration. Assuming both Houses have passed a bill, differences between their separate versions must be reconciled. Each House may accept or reject amendments made by the other House, or offer other amendments in lieu. If one House has rejected an amendment, the other House may nevertheless insist upon it. If a House insists upon an amendment that the other rejects, then the bill is lost unless the procedure set out in the Parliament Acts is complied with. Once a bill has passed by both Houses, or has been certified by the Speaker of the Commons as having passed the House of Commons in conformity with the Parliament Acts, the bill is finally submitted to the Sovereign for "Royal Assent". Since 1708, no Sovereign has failed to grant Royal Assent to a bill. Assent may be given by the Sovereign in person, but is usually given in the form of letters patent read out in each of the Houses; in the House of Lords the Clerk announces the Norman French formula "La Reyne le Veult", and the Bill thereupon becomes an Act of Parliament. In 1708 the formula used for the Scottish Militia Bill was "La Reyne s'avisera" (however, this was on ministerial advice). In theory the Sovereign has the right to either "withhold" or "reserve" the assent, however this right is not exercised. If assent were withheld, then the bill would fail. If assent were reserved, then formally a final decision on the bill has been put off until a later time; if Assent were not given before prorogation of the session, then the bill would fail. Private, Personal and Hybrid Bills. In the nineteenth century several hundred private Acts were passed each year, dealing with such matters as the alteration of local authority powers, the setting up or alteration of turnpike trusts, etc. A series of reforms has eliminated the necessity for much of this legislation, meaning that only a handful of private Acts are now passed each year. A private bill is initiated when an individual petitions Parliament for its passage. After the petition is received, it is officially gazetted so that other interested parties may support or contest it. Counter-petitions objecting to the passage of the bill may also be received. To be able to file such a petition, the bill must "directly and specially" affect the individual. If those supporting the bill disagree that such an effect exist, then the matter is resolved by the "Court of Referees", a group of senior Members of Parliament. The bill then proceeds through the same stages as public bills. Generally, no debate is held on the Floor during the Second Reading unless a Member of Parliament files a "blocking motion". It is possible for a party whose petition was denied by the Court of Referees to instead lobby a Member to object to the bill on the Floor. After the bill is read a second time, it is sent to one of two committees: the "Opposed Bill Committee" if there are petitions against the bill, or the "Unopposed Bill Committee" if there aren't. After taking evidence, the committee may return a finding of "Case Proved" or "Case Not Proved". In the latter case, the bill is considered rejected, but in the former case, amendments to the bill may be considered. After consideration, third reading and passage, the bill is sent to the other House, which follows the same procedure. If necessary, the bill may have to face two different Opposed Bill Committees. After differences between the Houses are resolved, the bill is submitted for Royal Assent. Personal bills relate to the "estate, property, status, or style" or other personal affairs of an individual. By convention, these bills are brought first in the House of Lords, where it is referred to a Personal Bill Committee before being read a "first" time. The Committee may make amendments or even reject the bill outright. If the bill is reported to the House, then it follows the same procedure as any other private bill, including going through an Unopposed or Opposed Bill Committee in both Houses. A special case involves bills that seek to enable marriages between those who are within a "prohibited degree of affinity or cosanguinity". In those cases, the bill is not discussed on the Floor and is sent at the committee stage to a Select Committee that includes the Chairman of Committees, a bishop and two lay members. Hybrid bills are public bills that have a special effect on a private interest. Prior to the second reading of any public bill, it must be submitted to the Clerk, who determines if any of the House's rules have been violated. If the Clerk finds that the bill does have such an effect on a private interest, then it is sent to the "Examiners", a body which then may report to the House that the bill does or does not affect private interests. If the latter, then it proceeds just like a public bill, but if the former, then it is treated as hybrid. The first and second readings are just as for public bills, but at the committee stage, if petitions have been filed against the bill, it is sent to a Select Committee, but the Committee does not have the same powers of rejection as Private Bill Committees. After the Committee reports, the bill is recommitted to another committee as if it were a public bill. Thereafter, the stages are the same as for a public bill, though, in the other chamber, the bill may have to be considered once more by a Select Committee. Supremacy of the House of Commons. Under the Parliament Acts of 1911 and 1949, the House of Commons is essentially the pre-eminent chamber in Parliament. If the Lords fail to pass a bill (by rejecting it outright, insisting on amendments disagreed to by the Commons, or by failing to vote on it), and the bill has been passed by the Commons in two consecutive sessions, then the bill may be presented for Royal Assent unless the House of Commons otherwise directs, and provided that the bill was introduced in the Lords at least one month before the end of each session. However, twelve months must have passed between the Second Reading in the first session, and the final vote on passage in the second one. Also, the bill passed by the Commons in each session must be identical, except to take into account the passage of time since the bill was first proposed. The effect of the procedure set out in the Parliament Acts is that the House of Lords may delay a bill for at least thirteen months, but would ultimately be unable to overturn the concerted will of the House of Commons. However, this procedure does not apply in the case of private or personal bills, nor to bills seeking to extend the life of Parliament beyond five years. Under the Parliament Acts, a special procedure applies to "money bills". A bill is considered a money bill if the Speaker certifies that it relates solely to national taxation or to the expenditure of public funds. The Speaker's decision is final and cannot be overturned. Following passage by the House of Commons, the bill can be considered by the House of Lords for not longer than one month. If the Lords have not passed the bill within that time, it is submitted for Royal Assent regardless. Any amendments made by the House of Lords are ignored unless accepted by the House of Commons. In addition to the Parliament Acts, tradition and conventions limit the House of Lords. It is the privilege of the House of Commons to levy taxes and authorise expenditure of public funds. The House of Lords cannot introduce bills to do either; furthermore, they are barred from amending supply bills (bills appropriating money to expenditure). In some cases, however, the House of Lords can circumvent the rule by inserting a "Privilege Amendment" into a bill they have originated. The Amendment reads: The House of Commons then amend the bill by removing the above clause. Therefore, the privilege of the Commons is not violated as they, not the Lords, have approved the tax or public expenditure. Delegated legislation. Many Acts of Parliament authorise the use of Statutory Instruments (SIs) as a more flexible method of setting out and amending the precise details for new arrangements, such as rules and regulations. This delegated power is given either to the Queen in Council, a Minister of the Crown, or to other named office holders. An Act may empower the Government to make a Statutory Instrument and lay it before both Houses, the SI to take legal effect if approved by a simple vote in each House; or in other cases, if neither House objects within a set time. In theory, Parliament does not lose control over such statutory instruments when delegating the power to make them, while being saved the necessity to debate and vote upon even quite trivial changes, unless members wish to raise objections. English Votes For English Laws. During the creation of the Devolved Administrations of Scotland and Wales, the idea of an English Parliament or Regional Assemblies were discussed but ultimately not implemented. This created an issue where the UK Parliament is acting as a "de facto" English Parliament on matters devolved to the national assemblies. MPs from all regions were free to debate and vote on issues which did not effect their constituencies or constituents. The Conservative Government of 2015 decided to address the issue in a controversial manner. Instead of bringing a bill to the Parliament, they proposed changes to the Statutory Instruments (SIs). Any bill brought before the Commons which is adjudged by the Speaker to only effect English Constituencies (or in some limited cases England and Wales) can have a ”double majority” rule imposed. In short, all MPs are allowed to debate and vote, but for a vote to be won both a count of votes of all MPs and a vote for English only MPs must be won. Privilege. Each House has a body of rights that it asserts, or which are conferred by statute, with the aim of being allowed to carry out its duties without interference. For example, members of both Houses have freedom of speech during parliamentary debates; what they have said cannot be questioned in any place outside Parliament, and so a speech made in Parliament cannot constitute slander. These rights are collectively referred to as "Parliamentary Privilege". Both Houses claim to determine their own privileges, and are acknowledged by the courts as having the authority to control their own proceedings, as well as to discipline members abusing the rules. Furthermore, each House is the sole judge of the qualifications of its members. Collectively, each House has the right of access to the Sovereign. Individually, members must be left free to attend Parliament. Therefore, the police are regularly ordered to maintain free access in the neighbouring streets, and members cannot be called on to serve on a jury or be subpoenaed as a witness while Parliament is in session. (Arrest for crime is still possible, but the relevant House must be notified of the same.) Parliament has the power to punish "contempt of Parliament", that is, violation of the privileges and rules of a House. Any decisions made in this regard are final and are cannot be appealed to any court. The usual modern penalty for contempt is a reprimand, or brief imprisonment in the precincts of the House, but historically large fines have been imposed. 

Haskell is a standardized functional programming language with non-strict semantics. Haskell features include support for recursive functions, datatypes, pattern matching, and list comprehensions. By combining these features, functions that would be difficult to write in a procedural programming language are almost trivial to implement in Haskell. The language is, as of 2011, the functional language on and in which the most research is being performed. The examples below give a taste of Haskell, oriented toward those familiar with other programming languages. Examples. The classic definition of the function: fac 0 = 1 fac n = n * fac (n - 1) A cute definition using Haskell's built-in list notation and the standard codice_1 function: fac n = product [1..n] Both definitions above should compile into the "same" efficient code via a smart compiler using equational reasoning. A naive implementation of a function which returns the nth number in the Fibonacci sequence: fib 0 = 0 fib 1 = 1 fib n = fib (n - 2) + fib (n - 1) A function which returns a list of the Fibonacci numbers in linear time: fibs = 0 : 1 : zipWith (+) fibs (tail fibs) The previous function defines an infinite list, which is possible because of lazy evaluation. One could implement codice_2 as: fib n = fibs !! n (codice_3 is an operator which gets the nth element of a list). The Quicksort algorithm can be expressed in Haskell concisely as: qsort [] = [] qsort (x:xs) = qsort [y | y&lt;-xs, y &lt; x ] ++ [x] ++  qsort [y | y&lt;-xs, y &gt;= x] Non-skipping stable merge sort is mgsort less [] = [] mgsort less xs = head $ until (null.tail) pairs [[x] | x &lt;- xs]  where  pairs (x:y:xs) = merge x y : pairs xs  pairs xs = xs  merge (x:xs) (y:ys)  | less y x = y : merge (x:xs) ys  | True = x : merge xs (y:ys)  merge xs ys = xs ++ ys where codice_4 is the function composition operator codice_5, codice_6 is a lambda expression (a nameless function), and the predefined function  codice_7  repeatedly applies a function "fun"  until a condition "cond"  is met; with codice_8 keyword introducing the local definitions showcasing the use of patterns (with codice_9 matching a non-empty list with the head codice_10 and tail codice_11) and guards. The Hamming numbers sequence is just hamm = 1 : map (2*) hamm `union`  (map (3*) hamm `union` map (5*) hamm) -- a union of two ordered lists: union (x:xs) (y:ys) = case (compare x y) of  LT -&gt; x : union xs (y:ys)  EQ -&gt; x : union xs ys  GT -&gt; y : union (x:xs) ys using sections (partially applied operators) and the built-in function codice_12 working with lists, whether finite or not, due to lazy (i.e. "by-need" ) evaluation. Also, enclosing function's name into back-quotes turns it into infix operator so that sub-expressions are automatically formed as if by placing parentheses appropriately, forming a nested expression, as in 2+3+5 --&gt; ((2+3)+5) ; and a comment line is introduced by two dashes. Finally, the infinite list of prime numbers by trial division is primes = 2 : sieve [3..] primes -- 2 : _Y ((3:) . sieve [5,7..])  where  sieve xs (p:ps) | (h,t) &lt;- span (&lt; p*p) xs =  h ++ sieve [n | n &lt;- t, rem n p &gt; 0] ps or as unbounded incremental Sieve of Eratosthenes with ordered lists, import Data.List.Ordered primes = 2 : _Y ((3:) . minus [5,7..] -- ps = [2..] \\ [[p*p, p*p+p..] | p &lt;- ps]  . unionAll  . map (\p -&gt; [p*p, p*p+2*p..])) _Y g = g (_Y g) -- = g (g (g (g (g ... )))) defining primes corecursively as natural numbers which are not multiples of primes. Or with arrays, import Data.Array import Data.List (tails, inits) ps = 2 : [n | (r:q:_, px) &lt;- (zip . tails . (2:) . map (^2)) ps (inits ps),  (n,True) &lt;- assocs (  accumArray (\_ _ -&gt; False) True (r+1,q-1)  [(m,()) | p &lt;- px,  let s = div (r+p) p * p, m &lt;- [s,s+p..q-1]] )] It's not just for quicksort! Haskell is also used to build "real world" programs, including [[../GUI/|graphical]] or web user interfaces. To get a taste of Haskell in the real world, check out web frameworks in Haskell and Haskell in industry articles on the Haskell wiki, or darcs, an advanced revision control system written in Haskell. Notes. [[/Not in book]] 

Proteins are a primary constituent of living things and one of the chief classes of molecules studied in biochemistry. Proteins provide most of the molecular machinery of cells. Many are enzymes or subunits of enzymes. Other proteins play structural or mechanical roles, such as those that form the struts and joints of the cytoskeleton. Each protein is linear polymers built of amino acids. Proteins are also nutrient sources for organisms that do not produce their own energy from sunlight and/or are unable to fix nitrogen. Proteins can interact with one another and with others molecules to form complexes. Index of chapters and main sections 

= Phonology = Examples. Feminine A=Stems (so is the fem. definite article) Weak Declensions. Weak nouns decline differently than strong nouns (hence the different name) and sometimes are more simple. Examples. But the Weak declensions are not all the same there exists a few irregulars. First the masculine irregulars: NEXT TO COME: FEMININE WEAK NOUNS, NEUTER AND FEMININE STRONG A-STEM VOCABULARY ALSO. = Links = Gothic language on the web 

Lektion 5 Wiederholung. Lesson 5 is a review ("Wiederholung") lesson to summarize the German language lessons presented in Lessons 1 through 4. You should, then, return to Lektion 1 and review (that is, reread) each of the four lessons back up to this point. For a more advanced course, you might now incorporate each of the advanced lessons into this "review" process. That is: review Lesson 1, then do Lesson 1A, review Lesson 2, then do Lesson 2A, etc. Parts of Speech and Word Order. Sentences are composed of parts that perform specific functions. You have been introduced to most (but not all) the major parts of speech: pronouns/nouns, verbs, and adjectives; and how these are expressed in German compared with English. Consider the following: Ich brauche Wurst und Käse Haben sie zu viel Arbeit? Word order in a simple sentence follows that used in English. Subject and verb are reversed to form a question. In English, but not in German, the question sentence could also be stated (and, in fact, occurs more often in the US) as 'Do they have too much work?' Nouns. Nouns are words that typically occur in sentences as either subjects (performers of some action) or objects (recipients of some action). Most nouns are the name of either a "person, place, or thing" and, in German, are always capitalized. Every noun in German has an "assigned" gender (masculine, feminine, neuter), and we learn each noun with its nominative case, definite article ("der", "die", "das", respectively) in order to also learn that gender. Thus, a "Vokabeln" section for nouns is presented thusly:  der Anhang, die Anhänge appendix, appendices (singular and plural)  die Brücke bridge  der Freund, die Freunde friend, friends (singular and plural)  das Gespräch, die Gespräche conversation, conversations  die Grammatik grammar (note irregular stress)  die Lektion lesson (note irregular stress)  die Straße street 

« Gothic Strong Nouns This is the strong noun A-Stem section, there are other kinds of strong nouns but we will deal with those later. Masculine A-Stem Nouns Gothic has strong and weak nouns (along with strong and weak adjectives!). There are many stems or different declensions within each class of nouns. Here I will show you a-stem nouns. There are three genders in Gothic: masculine, feminine, neuter. These are not always sex-based as in English. Now onto the masculine article (definite). Gothic has no indefinite article (a, an) but only a definite (the) which has masculine, feminine, and neuter forms. For a while here I'm going to be basically throwing you a lot of tables. But here I'll explain the cases: &lt;br&gt;Nominative: marks the subject of the sentence &lt;br&gt;Accusative: marks the direct object of the sentence &lt;br&gt;Genitive: marks possession &lt;br&gt;Dative: marks the indirect object of a sentence. 

Genealogy is the research of one's family history (or someone else's). What is included in "family history" varies depending on the individual researcher. There's a table of contents on every page that will help you navigate the pages. If you want to skip to a certain page, look below. Table of contents. Area-specific. /Authors/ Genealogy sites.  __NOEDITSECTION__ 

Strong Nouns This is the strong noun A-Stem section, there are other kinds of strong nouns but we will deal with those later. Masculine A-Stem Nouns Gothic has strong and weak nouns (along with strong and weak adjectives!). There are many stems or different declensions within each class of nouns. Here I will show you a-stem nouns. There are three genders in Gothic: masculine, feminine, neuter. These are not always sex-based as in English. Now onto the masculine article (definite). Gothic has no indefinite article (a, an) but only a definite (the) which has masculine, feminine, and neuter forms. For a while here I'm going to be basically throwing you a lot of tables. But here I'll explain the cases: &lt;br&gt;Nominative: marks the subject of the sentence &lt;br&gt;Accusative: marks the direct object of the sentence &lt;br&gt;Genitive: marks possession &lt;br&gt;Dative: marks the indirect object of a sentence. 

This book initiated on 02/14/04 by .&lt;BR&gt; Getting Started material by Diana Grzelak Needham.  is also working on this project.&lt;BR&gt;  from http://bob.fornal.org is also working on this project. 

Finding your family village name: You must find your ancestor’s community's name before you will be able to locate ancestral records that were recorded in the "old country." This will usually require beginning research in the U.S., or whatever your home country, to determine the correct community. First, begin with yourself. Gather everything you know about your own origins: where you were born, what religion you were raised in, your own civil and religious records. Then go back one generation at a time, looking for all of the available information you have. Where did your parents live? Where were they born? Gather their records. Sometimes even the smallest detail can lead you on the path to discovery. At times, you will not be able to locate a particular record for a parent or grandparent. If that is the case, look for records of their siblings. Often those will have the information you seek. =Records to Search= Family documents. Get out those old shoe boxes, start reading and really look at the details. Old letters often contain family connections, dates, locations and hints to lead you in the right direction. The same applies to old photographs which may have names and dates, or the photograpy studio or photographer’s name. Address books, postcards and, of course, family Bibles contain information and clues. Deeds, marriage certificates, naturalization papers, social security papers, Wills, and even insurance policies can be sources for information. I found one great-great grandmother’s address in Poland in an old personal address book which was about to be thrown out when her granddaughter passed on. Contact your relatives. Many families are large, and an aunt, great-uncle, or distant cousin may have family documents you were not even aware existed. Begin with the one who was closest to your ancestor or cared for them before they died. Often those relatives end up with the ancestor’s documents. An important note about family Bibles: in most of the older Bibles, the family information is listed about 2/3 to 3/4 of the way toward the back of the book (unlike modern Bibles where this information is listed in the first few pages). Public Records. Public records in general present some problems. Families changed their names; names were misspelled due to language barriers; phonetic spellings were used by clerks; different pronunciation of letters from one language to the next, typographical errors and poor penmanship created transcription errors, among others. Some records are missing due to misfiling, or were destroyed by fire and flood. This is as true in the U.S. as it is in other countries; the 1890 Federal Census being a prime example. Vital records. Birth, marriage and death records are generally located in the county seat in which the event happened (in the USA, other arrangements apply in the rest of the world). The cost of those records vary from county to county, and often depend upon whether you ask for a certified copy of the record or just a genealogical copy (generally cost less). Most U.S. counties can be accessed on the web as follows: www.co.”countyname”.”state two letter abbreviation”.us/ (omit the quotation marks in your search. If that doesn’t work, just do a web search with your favorite search engine. The information for obtaining the records should be on the county site. In the U.S. this information can be found on the USGenWeb site: http://www.usgenweb.org/. There is also a World GenWeb site that may help for researching in other countries: http://www.worldgenweb.org/. Another source to find vital records by state in the U.S. can be found at http://www.cdc.gov/nchs/howto/w2w/w2welcom.htm. Some of these state sites give the year in which the state began to maintain the records, the locations to find earlier records, and refer to databases which are also searchable on the web. SSDI. The Social Security Death Index is available free on line at http://ssdi.rootsweb.com/ Scroll down a bit and just put the name in the search box, and a list of persons with that name will come up. The Social Security Death Index is generally updated every three months. When a person is found, the name, date of birth, date of death, last residence, the state in which the number was issued, and the number itself is shown. If you wish to obtain a copy of the original Social Security Application, you can click on the SS-5 letter link. The following information was requested on the application: Name, current address, birthdate, birthplace (generally City and State, place of employment and business address (if applicable) and names of parents (mother’s full maiden name). Social Security is constantly updating this file. If the person for whom you are looking recently died, they probably won't be here. Give it a few months before you check again. Obituaries. Published obituaries can be a valuable source of information. Often parents’ names, spouses, children and surviving siblings are shown. This is especially valuable if you are unable to locate information on your ancestor and don’t have the female siblings’ married names. Generally, if listed at all, it is with their then present married name, giving you another piece of information to allow you to continue your search. Obituaries may be found in the “morgue” of the local newspapers, historical libraries in various cities and counties, and some public libraries have film or online databases for obituaries, marriages and death notices. Alien Registration. Alien registration was required beginning in 1940 for resident aliens in the USA. The records can be obtained for the years 1940-1944 from the Immigration and Naturalization Service. http://www.bcis.gov/graphics/aboutus/history/ImmRecs/AREG.htm Voter Registration. Most U.S. counties and states keep voter registration records. The states’ records are generally kept in the office of the Secretary of State. You can check with the county or state departments for information. Additionally, the LDS has filmed numerous voter registration indexes for various areas of the United States. You can check their Family History Library Catalog on line to see if your location has been filmed. http://www.familysearch.org/eng/Library/FHLC/frameset_fhlc.asp WWI Draft Registration. In 1917 and 1918 males old enough to serve in the military were required to register for the draft, whether citizen or alien. See more detailed information on Ancestry WWI Draft Information at Ancestry.com The LDS has compiled the WWI Draft Registration on microfiche arranged by state of residence, and you can contact your local Family History Center to view the registration records. Information is available at WWII Draft Info at FamilySearch.org United States Census. The U.S. census is an excellent source of information, although it won’t usually give you the town of origin, it will give the country and occasionally the Province, such as “Prussian Posen” or “Russian Warsawa”. Keep in mind that the reference to “Posen” and “Warsawa” in these records generally refer to the Province and not the city. Remember that these records are organized by family group. You will generally find parents and children, but sometime will also find grandparents and even great-grandparents living with their descendants. This information is also valuable because it provides a fairly accurate view of where the family (including siblings) were located when the individal census was taken - given several, this can show how the family members came and went in a household. Census films may be obtained at your local Family History Center of the LDS (see the Family History Library search page listed under voter registration above) or through the National Archives and Records Administration (NARA) http://www.archives.gov/research_room/genealogy/research_topics/census_records.html. These records are also available on line through the paid genealogical subscription sites, www.Ancestry.com and www.Genealogy.com. Immigration records. There are a number of sources for ship manifests, and their availability varies with the time and port to which your ancestor immigrated. The most well-known is the Ellis Island database on line, and if your ancestors immigrated between 1892 and 1924 to the port of New York, you may find that information free online: http://www.ellisisland.org/sign/index.asp?ACT=LL&amp;login_targ=none. You must register to use the site, but it is without charge. There have been some problems with the Ellis Island database indexing and linking, and the usual problems with misspellings and transcription errors. Another site which can help with such problems is the Stephen P. Morse website which gives search options not available from the Ellis Island site, and a way in which to find “missing manifests.” http://stevemorse.org/ For more information on Ellis Island records, go to the Jewish GenWeb site: http://www.jewishgen.org/infofiles/eidbfaq.html The LDS has filmed numerous passengers records for many ports, and have some ports indexed so that you can find your ancestor’s name and arrival date before you order the actual manifest. Again, you can check what records they hold at http://www.familysearch.org/eng/Library/FHLC/frameset_fhlc.asp NARA also has many manifests for U.S. ports: http://www.archives.gov/research_room/genealogy/immigrant_arrivals/passenger_records.html Additional passenger arrival information can be found on the following sites: Immigration and Ships Passenger Lists Research Guide http://home.att.net/~arnielang/ship04.html#whatrec Polish Genealogical Society of America http://www.pgsa.org/ships.htm An extensive reference for immigration to Australia, Canada, Europe, the U.S. and other parts of the world is found on Mary’s Genealogy Treasures at: http://www.telusplanet.net/public/mtoll/immigr.htm. Declarations of Intent or Citizenship Applications are generally found in the local county court of residence prior to the Federal government taking over the naturalization process. A guide to finding these records is found at: http://www.germanroots.com/naturalizationrecords.html There are some free and paid searchable naturalization databases on the web for the U.S. and Canada at: http://www.germanroots.com/naturalization.html Land records. These are often harder to obtain and study than other documents, but they may hold clues, if not outright references to former residence. Probate records. Again, these may hold clues, if not outright references to former residence. Religious records. Baptismal records and marriage applications kept by the churches often have parents’ names and place of birth. If you cannot find your ancestor’s batismal certificate, you may have be able to find a sibling’s record and locate the information indirectly. Catholic records are generally kept in the parish in which the ceremony was performed. If a parish has closed, the records may have been sent to the Diocese or Archdiocese archives, or to the nearest church still in existence. Many local parishes are now on line in the U.S. To locate a specific church by town or city and even by foreign Language masses, this site is very helpful: http://www.masstimes.org/ASP/ . There is even a scrollable map so you can find a parish nearest your ancestors’ home. There is no set web address for the Archidiocese in the U.S., but a web search for “catholic archdiocese” and the city or state will generally get you good results. For a Diocese, you can use a similar search. Many have information about how to obtain genealogical records. The Evangelical(Lutheran) church of the U.S. has genealogical assistance. Their website is: http://www.elca.org/os/archives/index.html. Genealogical information is at the top right. Additionally, the Lutheran World Federation has information and links to their member churches throughout the world. Just click on “member churches” and scroll down to the area in which you are searching: http://www.elca.org/os/archives/index.html. Jewish records: Jewish Temples did not (and do not) keep records as did the Christian churches. These are family records, and were kept in the family papers. Some of the mohels who performed the bris mila kept records, but they were his personal records and would be with his family documents. Of course, many of these were lost through war and the holocaust, but there are still many ways one can search for their Jewish ancestors. Jewis GenWeb has an excellent guide to searching for family and ancestors: http://www.jewishgen.org . 

Research success is more than knowing where to look. In addition to data about your family, you'll want to have maps, background information about places, material about resources and how to use them. It won't take long for you to gather so much information that you can't keep track of it all. A systematic method for handling data is crucial. A good system has these qualities: Decide how you will store irreplaceable materials (original documents, photos). Always copy or scan originals and store the originals in a secure location (safety deposit box, safe, firesafe box, etc.) Label each newly acquired document with a source description and an ID. You may find it easiest to just number each document sequentially. It becomes almost impossible to know in advance how to organize an ID system that handles all future acquisitions. Refer to this document ID when extracting data from the document into your research system (binders, software, etc.) Most genealogy software allows you to keep a list of sources. Use this ID to file the documents, or enter you own ID as part of each source notation. Unless you can always bring a notebook computer to every research site, you should use a binder to display the current state of your research for reference. You'll need to know what you already know as you look for new information. A new option for those that use computers for recording their research, you can use a USB thumb drive/flashdrive/memory stick and load the entire Personal Ancestral File (PAF) program on the drive along with your database and access it at any computer that has USB capabilities. Other computer programs may be able to do the same thing, but I am not sure. PAF is a free computer download from www.familysearch.org that has documentaion and multimedia abilities. The PAF program is a windows based program, but can be run on a MAC if a pc-emulator is used. 

When starting genealogy, you should ask yourself what you want to accomplish with your research. Are you primarily curious about your ancestors and want to see what information is easily available? Do you want to devote your life to uncovering and preserving the past? What parts of your family do you wish to research? Some people prefer to focus on a few lines and mostly ignore the lines that marry in. Some people collect every person that can be tied in, including every umpteenth cousin many times removed. What about adoptions? Which line do you research? If you have a surname to research that isn't very common, will you research anyone with that name, whether there is any known connection to you or not ("one-name study")? As time goes on, your scope may change. But try to determine at the beginning what you wish to research in order to avoid re-doing previous work. 

Working backwards from yourself to your most distant ancestors is a critical method in tracing one's own ancestry. Although it requires discipline, it is extremely important to systematically work backwards in time through the record trail, exhausting all record sources for each person, before working on their parents and grandparents. Without working systematically backwards, it becomes more and more easy to miss a crucial piece of evidence, or to rely on a hypothesis or guess as the foundation for later efforts. Use of a checklist for every person you are researching can make this methodical task easier. In order to properly work backwards, one should use evidence in one generation regarding the previous generation as your foundation. Birth, marriage, and death certificates, both civil and ecclesiastic, often name, or even describe the occupation of the bride or groom's father and mother. Many times the certificate's subject's place of birth is given, and sometimes the place of the parents birth as well. For example, old civil birth records (1865–1910) in New York City, found at the Municipal Archives, have a field for Father, Father's Place of Birth, Mother (with maiden name), and Mother's Place of Birth. Similarly, Old ecclesiastic records from the villages around Kaiserslautern in southern Germany often describe the fathers of both groom and bride, the fathers' residences, and occupations. Less often in the oldest records, the mother is named. However, using either record, one would certainly establish a solid footing for looking into the previous generation. However, a single record, such as a birth, is not enough if further records are available. If one readily find the civil birth record, it can serve as a guide for the location of the same person's ecclesiastic baptismal, christening, or circumcision records, which may confirm, conflict, or expand upon the information in the civil record. Further research should be conducted upon your subject, however. Marriage and death certificates and records, both civil and church, land records, coroner's reports, voters rolls, census returns, and numerous other pieces can add to your understanding of a person's parents and their history. In many cases, you will find information about several different generations at the same time. 

In this section we will look at certain mathematical processes which deal with the fundamental property of recursion at its core. What is Recursion? Recursion, simply put, is the process of describing an action in terms of itself. This may seem a bit strange to understand, but once it "clicks" it can be an extremely powerful way of expressing certain ideas. Let's look at some examples to make things clearer. Exponents. When we calculate an exponent, say "x"3, we multiply "x" by itself three times. If we have "x"5, we multiply "x" by itself five times. However, if we want a recursive definition of exponents, we need to define the action of taking exponents in terms of itself. So we note that "x"4 for example, is the same as "x"3 × "x". But what is "x"3? "x"3 is the same as "x"2 × "x". We can continue in this fashion up to "x"0=1. What can we say in general then? Recursively, with the fact that We need the second fact because the definitions fail to make sense if we continue with negative exponents, and we would continue indefinitely! Recursive definitions. Reducing the problem into same problem by smaller inputs. for example  a power n  2 power 4  the recursion(smaller inputs) of this function is = 2.2.2.2.1  for this we declare some recursive definitions  a=2  n=4  f(0)=1  f(1)=2  f(2)=2  f(3)=2  f(4)=2  for this recursion we form a formula f(n)= a.f(n-1)  by putting these smaller values we get the same answer. Recurrence relations. In mathematics, we can create recursive "functions", which depend on its previous values to create new ones. We often call these "recurrence relations". For example, we can have the function :"f"("x")=2"f"("x"-1), with "f"(1)=1 If we calculate some of "f"'s values, we get However, this sequence of numbers "should" look familiar to you! These values are the same as the function 2"x", with x = 0, 1, and so on. What we have done is found a "non-recursive" function with the same values as the "recursive" function. We call this "solving" the recurrence relation. Linear recurrence relations. We will look especially at a certain kind of recurrence relation, known as "linear". Here is an example of a linear recurrence relation: Instead of writing "f"("x"), we often use the notation "a"n to represent "a"("n"), but these notations are completely interchangeable. Note that a linear recurrence relation should always have stopping cases, otherwise we would not be able to calculate "f"(2), for example, since what would "f"(1) be if we did not define it? These stopping cases when we talk about linear recurrence relations are known as "initial conditions". In general, a linear recurrence relation is in the form The number "j" is important, and it is known as the "order" of the linear recurrence relation. Note we always need at least "j" initial conditions for the recurrence relation to make sense. Recall in the previous section we saw that we can find a nonrecursive function (a "solution") that will take on the same values as the recurrence relation itself. Let's see how we can solve some linear recurrence relations - we can do so in a very systematic way, but we need to establish a few theorems first. Solving linear recurrence relations. Sum of solutions. This theorem says that: This is true, since if we rearrange the recurrence to have "a""n"-"Aa""n"-1-"Ba""n"-2=0 And we know that "f"("n") and "g"("n") are solutions, so we have, on substituting into the recurrence If we substitute the sum "f"("n")+"g"("n") into the recurrence, we obtain On expanding out, we have But using the two facts we established first, this is the same as So "f"("n")+"g"("n") is indeed a solution to the recurrence. General solution. The next theorem states that: &lt;BR&gt;Then γ"r""n" is a solution to the recurrence, where "r" is a solution of the quadratic equation which we call the "characteristic equation". &lt;br&gt;We guess (it doesn't matter why, accept it for now) that γ"r""n" may be a solution. We can prove that this is a solution IF (and only if) it solves the characteristic equation ;&lt;br&gt; We substitute γ"r""n" ("r" not zero) into the recurrence to get then factor out by γ"r""n"-2, the term farthest on the right and we know that "r" isn't zero, so "r""n"-2 can never be zero. So "r"2-"Ar"-"B" must be zero, and so γ"r""n", with "r" a solution of "r"2-"Ar"-"B"=0, will indeed be a solution of the linear recurrence. Please note that we can easily generalize this fact to higher order linear recurrence relations. &lt;br&gt; Where did this come from? Why does it work (beyond a rote proof)? Here's a more intuitive (but less mathematically rigorous) explanation. &lt;br&gt; Solving the "characteristic equation" finds a function that satisfies the linear recurrence relation, and conveniently doesn't require the summation of all "n" terms to find the "n"th one.&lt;br&gt; We want : a function F("n") such that F("n") = "A" * F("n"-1) + "B" * F("n"-2) &lt;br&gt; We solve : "x"2 = "A"* "x" + "B", and call the solution(s) "r". There can be more than one value of "r", like in the example below! &lt;br&gt; We get : a function F("n") = "r""n" that satisfies F("n") = "A" * F("n"-1) + "B" * F("n"-2)&lt;br&gt; Let's check: Does "r""n" = "A"*"r""n"-1 + "B"*"r""n"-2 ? Divide both sides by "r""n"-2 and you get "r"2 = "A"*"r" + "B", which must be true because "r" is a solution to "x"2 = "A"* "x" + "B" &lt;br&gt; &lt;br&gt; Why does γ*"r""n" also satisfy the recurrence relation? If F("n")="r""n" is a solution, that is, "r""n"-"A"*"r""n"-1-"B"*"r""n"-2=0, then certainly F("n")=γ"r""n" is a solution since γ"r""n"-"A"*γ"r""n"-1-"B"*γ"r""n"-2=γ("r""n"-"A"*"r""n"-1-"B"*"r""n"-2)=0. &lt;br&gt;&lt;br&gt; Because we have a second order recurrence, the general solution is the sum of two solutions, corresponding to the two roots of the characteristic equation. Say these are r and s. The general solution is C(rn)+D(sn) where C,D are some constants. We find them using the two (there must be two so that we can find C and D) starting values of the relation. Substituting these into the general solution will give two equtions which we can (hopefully) solve. Example. Let's work through an example to see how we can use the above theorems to solve linear recurrence relations. Examine the function "a"("n") given here The characteristic equation of this recurrence relation is i.e. ("r"-2)("r"+1)=0 which has roots 2, -1. So the general solution is C(2n)+D(-1)n. To find C and D for this specific case, we need two starting values, let's say "a"("1") = 0 and "a"("2") = 2. These give a system of two equations&lt;br&gt; 0 = C(21)+D(-1)1 &lt;br&gt; 2 = C(22)+D(-1)2 &lt;br&gt; Solving these two equations yields: C = 1/3 , D = 2/3, so the solution is 1/3*(2n)+2/3*(-1)n. 

Presentation. With the breathing reality check, you check if you can breathe underwater or with your fingers tightly sealing your nose. Always remember to concentrate on the act. When you seal your nose, don't release the air pressure when you're done with the test. As well as doing this regularly, you could also do it every time you go swimming. NOTE: I have had personal experience with the breathing technique. I have used it numerous times. Only until recently it has worked perfectly fine for inducing lucid dreams; however, within the last two weeks, every time I try to plug my nose and breathe in, it doesn't work even in the dream world! It may be because I do the test so often that my expectation of the result carries into the dream world, therefore giving me the result I expect no matter if I'm dreaming or not. But, I do not see why this would happen since I constantly question the reality of the world around me hoping and expecting to be in a dream. My advice is to have a variety of techniques down in your mind and test two or three of them. Sometimes, my breathing AND hand checks don't work any more! I have found that text and digital clocks are the only solid way to know for sure since I seem to have exhausted the ability to use the techniques relating to my body. 

Presentation. Also called the flying reality check; with the jumping reality check, you jump up and see if you float back down. Also try these techniques: 

Recall the definition of a limit: A number formula_1 is the limit of a function formula_2 as formula_3 approaches formula_4 if and only if "for all" numbers formula_5 there exists a number formula_6 such that whenever In other words, given a number formula_9 we must construct a number formula_10 such that assuming we can prove moreover, this proof must work for "all" values of :formula_5 . Note: this definition is not constructive -- it does not tell you how to "find" the limit formula_1 , only how to check whether a particular value is indeed the limit. We use the informal definition of the limit, experience with similar problems, or theorems (L'Hopital's rule, for example), to determine the value, and then can prove the correctness of this value using the formal definition. Example 1: Suppose we want to find the limit of formula_15 as formula_3 approaches formula_17 . We know that the limit formula_1 is 9+5=14, and desire to prove this. We choose formula_19 (this will be explained later). Then, since we assume we can show which is what we wanted to prove. We chose δ by working backwards from the formula we are trying to prove: formula_7 . In this case, we desire to prove given so the easiest way to prove it is by choosing formula_19 . This example, however, is too easy to adequately explain how to choose formula_10 in general. Lets try something harder: Example 2: Prove that the limit of formula_26 as formula_3 approaches formula_28 is formula_29 . We want to prove that given We choose formula_10 by working backwards. First, we need to rewrite the equation we want to prove using formula_10 instead of formula_3: formula_35 Note: we used the fact that formula_36 , which can be proven with the triangle inequality. Word of caution: the above series of equations is not a logical series of steps, and is not part of any proof, but is an informal technique used to help write the proof. We will select a value of formula_10 so that the last equation is true, and then use the last equation to prove the equations above it in turn (which is what was meant earlier by "working backwards"). Note: in the equations above, when formula_10 was substituted for formula_3 , the sign formula_40 was replaced with formula_41 . This can be done (but is not necessary) because we are not told that formula_42 , but rather formula_31 . The justification for this becomes clear when the above equations are used in backwards order in the proof. We can solve this last equation for formula_10 using the quadratic formula: Note: formula_10 is "almost always" in terms of formula_9 . A constant value of formula_10 (e.g., formula_49) will not work unless the limit is of a constant function (for instance, formula_50). Now, we have a value of formula_10 , and we can do our proof: given formula_53 Here a few more examples of choosing formula_10 ; try to figure them out before reading the explanation. Example 3: Prove that the limit of formula_55 as formula_3 approaches formula_57 is formula_58 . Explanation: Example 4: Prove that formula_59 has no limit as formula_3 approaches formula_57 . Example 5: Prove that formula_62 Solution: To do it, we'll look at two cases: formula_63 and formula_64 . The formula_63 case is easy. First let's let formula_66 . That means we want the values chosen in the domain to map to formula_67 in the range. We want a delta such that formula_68 so let's choose formula_69 . The chosen formula_10 defines the interval formula_71 in our domain. This gets mapped to formula_72 in our range, which is contained in formula_67 . Notice that formula_10 doesn't depend on formula_9 . So for formula_76, we widen the interval in the range that we are allowed to map onto, but our interval in the domain stays fixed and always maps to the same sub-interval in the range. So formula_69 works for any formula_63. Now suppose formula_79 . We want a formula_10 such that formula_81 whenever formula_82 . So let's assume formula_81 and work backwards to find a suitable formula_10: Since formula_79 , we have formula_89 . Since both numbers above are positive, we can take the (positive) square root of both extremes of the inequality: The above equation represents the distance, either negative or positive, that formula_3 can vary from 2 and still be within formula_9 of 4. We want to choose the smaller of the two extremes to construct our interval. It turns out that formula_94 for formula_79 , so choose formula_96 . As a sanity check, let's try with formula_97 . which is approximately At the extreme right of the domain, this gives and which is within 0.002 of 4. 

 This book aims to teach you how to use the vi  editor, common to many Unix and Unix-like operating systems.  ~  ~  "Learning_the_vi_editor" [New file]. The above text is a little example of how the vi editor's screen looks. 

Introduction. Overview. vi is a powerful editor that is ubiquitous amongst Unix and Unix-like operating systems, but is available on many other operating systems, even on MS-DOS, Windows and the Macintosh. If not the original vi, there is usually at least a good clone available that runs on your system. Even if you use another editor you must have a passing knowledge of vi as an administrator. Sometimes vi is the only editor available when your computer crashes leaving a minimal system for you to repair. vi, pronounced like 'vee eye', was originally written by Bill Joy for BSD Unix in Berkeley in 1976 and became quickly part of many vendor-specific versions of the (at that time) original AT&amp;T Unix. It was later directly added to AT&amp;T's System V Unix, too. Bill Joy later went on to co-found Sun Microsystems, and became the company's Chief Scientist at that time. vi stands for "visual" and was an enormous improvement of the classic Unix editor called ed. ed is a line-editor. If you are still familiar with MS-DOS, then you may know the MS-DOS edlin editor. ed is similar, although more powerful than edlin, which doesn't mean much. vi also has a line-mode, called ex. In fact, one can argue that the program is indeed two editors in one, one editor called vi, another called ex. It is possible to switch between line and visual mode during editing. It is also possible to choose the mode during startup. However, pure usage of ex is rare. The visual mode is the prevailing mode. Although vi stands for "visual", classic vi is mainly operated via the character keys, and not via the mouse or the cursor keys. Once you are used to this, it becomes extremely convenient, because there is less movement of the hands to the cursor keys or mouse involved. vi also served as a kind of incubator for Unix's terminal control capabilities. Because of vi's need to control the terminal and the many different types of terminals at that time, the "termcap" (terminal-capabilities) database was introduced (later replaced with the more flexible "terminfo" database). vi's internal high-level screen control library was later separated, and became "curses" - the Unix standard library for CRT screen handling. Conventions. unix-command(section) Getting vi if you don't have it already. If you're running a Unix system, or a Unix-like system (for simplicity from now on we will refer to both as a "Unix system"), such as a BSD or Linux distribution, or even Mac OS X, you're sure to have vi or one of its variants on your system. If you're running Windows, you can get a version of vi called "vim" or "elvis". If you're on an older Mac OS (pre-OS X) system, you can get MacVim Classic here. Noted vi variants. As mentioned, vi has a number of variants. They have been created because vi was only available on rather expensive Unix operating systems. Although vi itself, as well as nvi, was created in Berkeley for the free BSD Unix variant, usage of BSD Unix required an original AT&amp;T Unix license (this has later changed, see below). Original vi, for example, used code from AT&amp;T's ed"(1)" editor. Over time, BSD replaced many of the original AT&amp;T code up to the point where today there is no such code any more in BSD, and an original Unix license is no longer needed. As part of the effort to replace all AT&amp;T code in BSD, Keith Bostic undertook the work to create a clone of vi that was free of AT&amp;T code, called nvi. nvi then became BSD's standard vi instead of the original vi. Another vi clone is Elvis, which was written by Steve Kirkendal. Over time, nvi was enhanced – for example, supporting multiple windows – but originally it was not supposed to be an enhancement, 'just' a pure clone. BSD's original vi (with the ed code inside) lives on as the vi which is distributed with System V Unix, because AT&amp;T decided a long time ago to take it from BSD and add it to the official Unix. Of course AT&amp;T didn't have a problem with an AT&amp;T Unix license, so they probably never replaced the ed code inside the original vi. Yet, some find nvi still to be too minimal, and so vim was born. vim (vi-i"m"proved), is another effort to extend vi's capabilities. Unlike nvi, vim goes even further to extend vi's capabilities. However some find that vim is often too much. vim comes in two variants, a text-only version, and a GUI version, the latter is called gvim. Other vi clones are the already mentioned elvis and stevie. These clones were born in the CP/M and home computer area to bring the editor to these platforms, too. Of course, they were later ported to MS-DOS and Windows. These days, however, vim seems to be the prevailing vi-clone on free/open platforms and proprietary platforms as well. "You should choose the version you feel most comfortable with" – if you have an editor you feel displeased with, it will affect your productivity. Getting around vi. Starting the editor. If you are running a Unix system, you can start up vi by typing  vi at the command line. If you are running X, with a desktop like GNOME, KDE, CDE/Motif or OpenLook you may have a launcher button handy to start the editor - if you have such a setup, you can just click the icon. If you are running Windows or DOS with elvis, you can start up the Windows editor by double-clicking "winelvis.exe", or in DOS, you can start the editor by typing in "elvis" at the command line. You will be greeted with a screen similar to:  ~  ~  ~  "No File" Quitting the editor. To quit for now, press the Escape key (the editor should beep), then enter the three characters and press Return: Just before you type the final the screen will look similar to  ~  ~  ~  :q!  is the short form of which quits the editor. You should be dropped back to your operating system (or, rather, the shell from where you started). There are other ways to quit, e.g. pressing () will save any unsaved work and quit the editor. Typing will always save, even if there are no unsaved changes, and then quit the editor. will write if there are no unsaved changes, and it will quit. and requires that you had previously provided a file name, so it will not work for the above simple example. Typing will quit if there have been no changes made; if changes have been made, vi will print a warning similar to "No write since last change". Don't worry. Many first time vi users stop at this point, and never touch vi again. If you tried to enter some text after you started, you will most likely have been greeted with a series of beeps and rather erratic behavior. Don't worry. This is perfectly normal for vi, and the editor is not broken. You will soon see why this is normal vi behaviour. Continue. Now that you know how to start the editor and quit it, let's move on to getting things done in vi: see Learning the vi Editor/Basic tasks 

Now that we know how to invoke the editor and quit it, we can get acquainted with how to "use" the editor. Alternatively, you can use the ViM tutor which comes with many modern vim distributions. It contains, essentially the same information as the text below. You can invoke the tutor by entering vimtutor at your shell. Vi is a "modal" editor. The vi editor can do two things: In the vi editor, each of these tasks is achieved by putting the editor into a particular mode of operation (normally just called a "mode"). When you wish to give vi a command, you enter "command mode", and when you want to enter text, you enter "insert mode". We'll cover how to do this below. It is important to set the correct mode before you begin writing, but this is simple to do. When you first start vi, it is automatically in command mode. Entering text. Entering text is the most basic task an editor can do! From command mode (in which the editor starts), press codice_1 to enter "insert mode", and you can begin typing. You can use the backspace key to correct mistakes you make. If you make a mistake after a few sentences, leave these errors for now, we will look at correcting them later. To leave insert mode once you're done typing, and return to command mode, press the Escape key on your keyboard (or type Control-[). Exercise. Let's have an exercise: Command mode. Command mode allows you to perform many useful tasks within vi. Moving around. Say you have been writing for some time, and have forgotten something. Pressing , "erasing" previous work is not the best solution! We would like to move around the document freely, moving the cursor. Exercise. You can repeat this exercise with your own sentences. Make sure you are proficient doing this before you continue. More on movement. Using codice_3, codice_4, codice_5, and codice_6 is ok, but vi understands more than rows and columns. These are some commands that move by text objects: Deleting things. If you have made a mistake after a few lines, for instance, pressing Backspace until you have erased the mistake and starting again isn't always the best solution. We need a method of deleting mistakes that happen in the normal course of editing. vi allows you several methods of deleting text, based on how much you want to remove. Now that you are familiar with moving around, once you've moved the cursor to where your error is: Exercise. From now on, we will omit the steps for you to start and quit the editor – you should be already familiar with those. 



The Mathematics Of Waves. We start our discussion of waves by taking the equation for a very simple wave and describing its characteristics. The basic equation for such a wave is where formula_2 is the height of the wave at position formula_3 and time formula_4. This equation describes a fairly simple wave, but most complex waves are just sums of simpler ones. If we freeze this equation in time at formula_5, we get which looks like this: [TODO - Add a Graph] From the graph we can see that each of the three parameters has a meaning. formula_7 is the amplitude of the wave, how high it is. formula_8 is the wavelength, the distance from a part of the wave in one cycle to the same part of the wave in the next cycle. formula_9 is the phase of the wave, which shifts the wave to the left or right. The wavelength is a distance, and is usually measured in meters, millimeters or even nanometers depending on the wave. Phase is an angle, measured in radians. Now that we have mapped out the wave in space, let's instead set formula_10 and see how the wave changes over time Amplitude formula_7 and phase formula_9 remain, but the wavelength is gone and a new quantity has appeared: formula_14, which is the frequency, or how rapidly the wave moves up and down. Frequency is measured in units of inverse time: in a fixed period of time, how many times does the wave move up and down? The unit usually used for this is the hertz, or inverse second. Now let's combine these two pictures and see how the wave moves. Figure 3 is a diagram of how the wave looks when you plot it in both space and time. The straight lines are the places where the simple wave reaches a maximum, minimum, or zero (where it crosses the x axis). We can look at the zeros to determine the phase velocity of the wave. The phase velocity is how fast a part of the wave moves. We can think of it as the speed of the wave, but for more complicated waves it is only one type of speed - more on that in later sections. We can get an equation for the zeros by setting our equation to zero. You see here that we have the equation for a straight line, describing a point that is moving at velocity formula_18. This gives us the equation for the phase velocity of the wave, which is 

Appendices. The following are Java code referenced in the book. 

Code Generation. A compiler usually is designed to output an executable program that will allow the user to run your program, and to be directly run by the processor, without having an intermediary interpreter such as in the interpretation process. For your program to be run by the processor however, you will need to transform the instructions in your specific programming language into assembler code, which is then sent to an assembler tool to create object code, which is then linked together with specific libraries to create your executable code. For now, we only really need to worry about transforming the instructions into assembler code. This process is what we will deal with in this section. You will need to be well versed in the assembler language you wish to output. If you intend your programs to run on the x86 architecture, you need to be familiar with x86 assembler code, and so on. Code generation occurs after semantic analysis is done, which gives us enough information to generate more primitive concrete code. The general idea behind code generation is decompose the tree structure of the syntax tree into a sequence of instructions, whatever an instruction set is. In this stage, since we are done with the semantic program, we are not interested in the syntactic and semantic structure of programs but in the order of executions of instructions. Some sort of intermediate code is often produced before generating actual machine code. The benefits of this are However, it may be simpler for your program to output assembler code directly, but you lose the above advantages. See the next section for more techniques on this. In this chapter, we shall use the three address format to represent intermediate code. The format is useful because it is analogous to actual machine instructions in some architectures and, more importantly, allows us to easily change the execution order of instructions, which is a huge advantage over stack-based intermediate code like the byte code of Java. Although it is not a complex problem to reuse names after they have already been used, it is actually beneficial to allocate a new name every time one is needed because it allows us to form a call graph and optimize easily as we will see later. For this reason, we only briefly mention the methods to reuse names. You can find more on the optimization of allocation of names in optimization chapter. The three address code, as the name suggests, consist of three address and opcode, which tells what kind of operation is meant to be done. For example, an expression (a + b) * 3 can be transformed into: temp1 := a + b; temp2 := temp1 * 3 In the first line, temp1, a and b are addresses and + is an opcode, and the second line is similar to the first one. Unlike load-store machines, it is unnecessary to load variables to registers and store them back. You see why the three address code is easy to handle. Choosing portable, flexible and expressive instructions is critical; Not having enough instructions can complicate generated code with the combination of several instructions to achieve one operation and having too much may obviously make maintenance more daunting task. Probably the best way to do this is to examine existing machine code. It is more straightforward to transform code close to underlying machine code than abstract one. Expression. Algebraic expressions can be translated into the three address code in a very straightforward manner. This can be done rather recursively as follows: Assume two expressions left and right with an operation op-code, then the results should be:  code for left  code for right  temp = place for left + place for right Control Structures. The general idea behind generating code for control structures is the same as coding in assembly programming. That is, an if statement, for instance, is converted into a chunk of code using conditional and unconditional jumps. Assembler language techniques. If the output of your compiler is assembly language code, it is necessary to understand the basic techniques of assembly language programming. Most programming languages do not map easily to most assembler languages, so some techniques or skills may need to be understood before attempting to write code that will output assembler code. These techniques are not intended to create highly optimized code - you will learn optimizing techniques later - but are intended to make sure you have a good understanding of how data and instructions are managed in the process of compiler construction. Managing variables. Many programs use hundreds of different variables (not counting arrays). Most computer architectures give you less than 32 registers (MIPS architecture and ARM give nearly 32 pointer registers; i386 gives only about 4 pointer registers; PIC microcontroller only has 1 pointer register). Since we can't squeeze 100 different variables into even 32 processor registers, we must use memory for storing most variables. We will start by storing practically all variables in memory. Later we will cover optimizing techniques that try to keep as many variables as possible in the processor registers. Labeling. The assembler tool that you are using may reserve the names of the mnemonics. For example, your assembler may not allow a variable named add, since this is reserved for the instruction to add. In this case, it may be important to use a "prefix" for your variable labels. Some compilers use a single underscore, but you can choose whichever you wish. 

Introduction. This textbook concerns the wonderful world of macroeconomics, or economics on a very large scale, concerning national and international systems. It is primarily aimed at students in their final few years of secondary education, though it could also be used by interested students younger or older than that. Having a background knowledge of would be useful, but grasping overall concepts should not require an in depth previous knowledge. It is worth remembering that this textbook can be edited at any time, with the link at the top of this page. This is both good and bad - you yourself, having spotted a mistake or having noticed a poor definition, can correct it and should feel free to do so. The bad side is that anyone can edit it, so content may be inaccurate whilst the wikibook is in its infancy! Proposed Change. Advanced topics Further chapters later Useful Background Reading. The following links will lead you to Wikipedia, the online encyclopedia similar in style to this site. The information in these articles may go above and beyond that which you need to know for this book, though they will undoubtedly prove useful to refer back to from time to time, and are good portals to more relevant articles. The following links take you to useful websites for learning, revising or expanding your knowledge: For contributors. Any and all corrections, additions, alterations etc. are always welcome! If you wish some credit for your work, stick your name below : 

Welcome to the web's free, open content textbook of = Personality Type and the Myers-Briggs Type Indicator = 

The Myers-Briggs Type Indicator (MBTI) is the name of a personality test designed to assess psychological type. It was developed by Katherine Briggs and her daughter Isabel Myers during World War II. The use of type follows from the theories of Carl Jung. The phrase is also sometimes used as a trademark of CPP Inc., formerly known as Consulting Psychologists Press, Inc. The trademark is owned by the Myers Briggs Type Indicator Trust, and when used as a trademark it must include a registered trademark symbol after the name: "Myers-Briggs Type Indicator®" or "MBTI®". There are a few widely used ways of interpreting the results: Jung-like methods, Keirsey-like methods, and popular psychology methods. The MBTI is popular with recruiters and managers, because studies using this assessment show clusters of different personality types in different professions. For instance, the proportion of engineers who are ../INTJ/ is higher than the 1% found in the general population. There are significant differences by sex, especially on the T vs. F distribution. Proponents of the system claim that almost all arguments between people tend to be manifestations of a type conflict (e.g., E vs. I, S vs. N, T vs. F, J vs. P). The P-J conflict is said to be the clearest: one person gets mad when the rules are broken and the other gets mad when rules are made. The Myers-Briggs Type Indicator is perhaps the world's most popular personality type description tool. Since its inception, many people have turned to the MBTI® for a deeper understanding of themselves and others. This MBTI textbook is designed to bring together general knowledge about the Myers-Briggs Type Indicator and make it available for understanding and application for individuals and groups, for personal and professional lives. 

The Myers-Briggs test describes four basic areas of personality: It is important to keep in mind that each dimension reveals a person's inborn preference with how he or she is most comfortable operating, and does not say that any person will "always" retain his preferred dimension. 

There are sixteen "types": 

There are four temperaments: SJ, SP, NT, and NF. SJ. Sensing judgers, or traditionalists, are practical people that keep the home fires burning and businesses working. They're always aware of who owns what and which social positions are held by whom. Their quest is to run everything, and they are often good at doing so. SP. Sensing perceivers, or experiencers, are adventurous, fun-loving, observant, physically skillful, impatient, easily bored, and good with tools and art. They want to be happy and make others happy too. NT. Intuitive thinkers, or conceptualizers, are analytical, impersonal, intellectual, rather unworldly, absent-minded, and are more likely to forget appointments. They continually seek to acquire new skills and pride themselves on their skills, logic, and efficiency. NF. Intuitive feelers, or idealists, understand people and tend to be aware of people's feelings. They can be warm, sympathetic friends, but find offense in the smallest careless remark. They tend to be very skillful negotiators and good with words. 

The MBTI is not yet scientifically proven. Skeptics, including many psychologists, argue that the MBTI has not been validated by double-blind tests (in which participants accept reports written for other participants, and are asked whether or not the report suits them) and thus does not qualify as a scientific assessment. Some even demonstrate that profiles can apparently seem to fit any person by confirmation bias, ambiguity of basic terms. and the Byzantine complexity that allows any kind of behavior to fit any personality type. See this extensive skeptical treatment of the subject. A Temptation to Pigeonhole. Another argument says that, while the MBTI is useful in self-understanding, it is commonly used to pigeonhole people or for self-pigeonholing. Supporting arguments include: 



The Introverted Sensor Thinker Judger is a traditionally masculine personality type, detached, unexpressive, concrete, and not given to speculation. Working with an ISTJ. ISTJ's are detail people. You will mostly find them in technical service industries, including medical fields, accounting, police work, engineering, office managers, and so forth. They are keen sensors of their environments; tastes, smells, noise levels. They often prefer a quiet work environment, where they can concentrate and focus on the task at hand. They work slowly and methodically, and don't like to be rushed. They can be focused on keeping structure in order and feel most comfortable when everyone is doing what they are supposed to be doing. ISTJ's value their personal privacy, and like to remain professional. Due to their discreet and reserved nature, others may see them as snobbish or uptight. Often times they remain anonymous for achievements they have accomplished. They may be upset with people who don't follow the rules, honor their promises, or follow through, but may not say so publicly. Traits they may value among fellow coworkers are loyalty, friendliness, and a strong work ethic. They are usually willing to lend their help to others in need if they have time. ISTJs in relationships. ISTJ Relationships in the general sense are regarded as of high importance, and are expected to be of same value to associates. 

ENFPs are outgoing and talkative types who often have many diverse friends. They are naturally irreverent, curious, good with words, creative, and even a bit artsy. They favor abstraction over detail and tend to flit from one activity to the next, deriving more pleasure from starting projects than from finishing them. ENFPs at work. ENFPs are sponges of knowledge. Curiously collecting information about life's experiences as fodder for their own stories, they live their lives much like a dramatic character in a movie, delighting in finding irony and drama in life situations wherever they go. They interpret the data they gather from people and the world in the form of inner meanings, relationships, and possibilities, often recounting their ideas to people though conversation. They love to spread the word! ENFPs, ../NF/s, are interested in many different things and this often leads them in multiple directions, therefore they are known for changing jobs and even careers frequently. ENFPs prefer to make work fun and have been known to put off work when there is a chance to have fun. ENFPs prefer environments which have little administrative or detailed work. Delegation to reliable support staff is important for success. ENFPs seek out jobs not confined by strict schedules and mundane tasks. According to Keirsey, the "Idealists" want to make the world a better place to live by helping others and are on a search for their "true self". Integrity, ideals, beliefs, and values matter to ENFPs. ENFP's practice and dream about being the best at what ever they seek to achieve. They want to KNOW everything! ENFPs generally have the following traits:  * Project-oriented  * Bright and capable  * Warmly, genuinely interested in people; great people skills  * Extremely intuitive and perceptive about people  * Able to relate to people on their own level  * Service-oriented; likely to put the needs of others above their own  * Future-oriented  * Dislike performing routine tasks  * Need approval and appreciation from others  * Cooperative and friendly  * Creative and energetic  * Well-developed verbal and written communication skills  * Natural leaders, but do not like to control people  * Resist being controlled by others  * Can work logically and rationally - use their intuition to understand the goal and work backwards towards it  * Usually able to grasp difficult concepts and theories Common careers for ENFPs are those which allow flexibility and spontaneity while serving others. This includes teaching, general practice of medicine, medical research, religious fields, entrepreneurship, missionary work, social work, community development, creative arts, acting, broadcasting, consulting, coaching, corporate training, public relations, counseling, advertising, and marketing. They enjoy working for themselves and creating new ideas to help people. The perfect job for an ENFP is one where they would give out useful information to others and they would receive several hundred dollars per hour for that service. ENFPs in relationships. Strengths ENFPs are energized by being around people. ENFPs take their relationships very seriously. ENFPs seek and demand authenticity and depth in their personal relationships, and will put forth a lot of effort into making things work out. ENFPs are generally warm, considerate, affirming, nurturing, and highly invested in the health of the relationship. ENFPs have excellent interpersonal skills, and are able to inspire and motivate others to be the best that they can be. ENFPs are generally highly valued for their genuine warmth and high ideals. ENFPs want to help, to be liked, and to be admired by other people on both an individual and a humanitarian level. ENFPs hold up their end of relationships, sometimes being victimized by less caring individuals. Weaknesses ENFPs are well aware of their weaknesses, often being their own harshest critic. Often the ENFP is reluctant to share intimate feelings unless in the company of deeply trusted relationships. Despite being extroverts, ENFPs require a great deal of alone time to center themselves. ENFPs who don’t have a positive support system may be strongly influenced by the opinions of others. ENFPs can exhibit preoccupation in their relationships, sometimes putting "all their eggs in one basket" and can tend to hyper focus on the other individual, in attempts to "fix" the other person or pull out their "real" emotions, transforming them into the perfect person. ENFP who do this are doing so because they don’t want to focus on their own interpersonal issues at hand. In this scenario, ENFPs may feel very anxious and preoccupied if the other partner is silent, non expressive, or withdrawn when coping with stress, instead of talking through things. This can deeply hurt them. Although energetic and effervescent, the ENFP can sometimes be smothering in their enthusiasm. They do not understand why someone would not be charmed by their enthusiastic display of affection and quirky jokes, because many are naturally drawn to their personality. They may try too hard to be what others want them to be, showing codependent tendencies. ENFPs are seen as disorganized by OCD types, having a short attention span, taking on too many tasks at once, lacking structure, and all over the place. Others expect the ENFP to follow societal roles and ENFP are a role unto their own. How to Behave Toward an ENFP. Give ENFPs the freedom to be flexible. Realize that churning through possibilities inspires their minds and get their creative juices flowing. Don't bog the ENFP down with too many details, especially on any subject not known to be of deep interest to him or her. Involve the ENFP in the process and try to keep things fun. When communicating with ENFPs you will find that many like to talk. Do not hesitate to interrupt and state your opinion. ENFPs enjoy speculating about ideas. They enjoy others who engage them in conversation, contribute their ideas and keep it moving in a positive direction. Overall ENFPs appreciate honesty in others, they want to know how people really feel. ENFPs are easily influenced by what other people say. It may not appear that they are listening when they are talking to you, but soon after the conversation ends they often ponder what the other person has said and incorporate those ideas into their own thinking to use the knowledge for the future. The main points to remember, ENFPs: &lt;br&gt; 1) Want to help and please the people they are working with, so give them frequent feedback.&lt;br&gt; 2) Like to hear from their significant others that they are loved and valued, (and are willing and eager to return the favor,) so let them know what you appreciate about them. &lt;br&gt; 3) Prefer happy and upbeat relationships, nevertheless, when conflict occurs, they usually want to engage in a dialog to work it out.&lt;br&gt; 4) Are one of the most freedom orientated personality types. Give them room to have some adventures on their own.&lt;br&gt; 5) May have a difficult time staying focused and following things through to completion. Let them know clearly what your deadlines are. ENFPs want to please other people, so let them know how important it is to you and why it is important to you.&lt;br&gt; 6) Prefer a participative and collegial atmosphere in which employees are included in the decision making. 

= NF (The iNtuitive Feeler) = NFs are categorized as idealists and share the traits of dominant iNtuitive and Feeling traits. They include the ../ENFJ/, ../ENFP/, ../INFJ/, and ../INFP/. The iNtuitive side of the NF makes him or her more inclined towards the abstract and the future as opposed to the concrete present. The Feeler preference means he makes decisions based on his personal values more than on cold analytical data. 

Introverts are rested and energized by solitude, and are very effective in solitary pursuits. An introvert is a person who prefers to process thoughts internally. Introverts tend to think before they speak. The word is also used informally to refer to somebody who prefers solitary activities to social ones, which is more of a behavioural than cognitive definition. Introverts tend to be seen as quiet and reserved, which is often confused with a lack of confidence by louder, more extroverted people. They often perform well in analytical roles that require intelligence or logic, but place less emphasis on social interactions and "people skills". 

Extraverts are energized by interacting with other people. They can often appear to be outgoing and may be effective in pursuits that involve interacting with other people. Extraverts tend to be sensation-seeking, spontaneous, and gregarious. They may enjoy crowds, noise, and stimulation. In a conversation, an extravert will tend to talk faster and louder than an introvert and to interrupt more frequently. 

Sensors want, trust, and remember facts, and usually describe themselves as "practical". For a Sensor, intuition is untrustworthy, and might seem like mental static. Sensation, as a perceiving mode of consciousness, focuses on heightening reality. Guardians share the combination SJ, while Artisans share the combination SP. 

Intuitives prefer metaphor, analogy, and logic, and tend to reason from first principles and hunches. While Sensors pride themselves on living in the real world, Intuitives pride themselves on seeing possibilities. This can cause conflict. Intuition, as a perceiving mode of consciousness, filters experience through the unconscious mind. Intuition focuses on possibilities rather than realities. Idealists share the combination NF, while Rationals share the combination NT. 

Thinkers use impersonal means of reasoning: logic and verifiable experience. 

Feelers prefer personal reasoning: value judgements and emotions. Thinkers often find Feelers muddle-headed. Feelers often find Thinkers cold and inhuman. 

Judgers prefer to come to decisions and move on. They can feel betrayed if a decision is "reopened". They are prone to hastiness, but get things done. 

Perceivers prefer to leave their options open to perceive new possibilities and processes as long as possible. They tend to mourn opportunities lost to premature decisions. They are prone to analysis paralysis, but rarely make permanent mistakes. 

The terms group theory and ring theory are refinements of algebraic understanding that developed in the era of electronics and aircraft, the 20th century. The term hypercomplex number harkens back to the age of steam. For the most part, the hypercomplex systems have been assimilated through the resolution of vision provided by groups, rings, and fields, and the term has been retired from use other than historic reference. Similarly, the field of complex numbers formula_1 has an insufficiently descriptive name, and might be better described as division binarions C according to composition algebra theory. Hypercomplex numbers grew out of William Rowan Hamilton's construction of "quaternions" in the 1840s. The legacy of his vision continues in spatial vector algebra: for vectors formula_2 and formula_3 the well-known products are These products are the severed remnants of Hamilton’s quaternion product: formula_6 In 1845 John T. Graves and Arthur Cayley described an eight-dimensional hypercomplex system now referred to as "octonions" or "Cayley numbers". They extend quaternions but associativity of multiplication is lost. James Cockle challenged the presumption of quaternions in four dimensions by presenting associative hypercomplex systems "tessarines" (1848) and "coquaternions" (1849). Hamilton had his own eight-dimensional system ("biquaternions") that were explored in his "Lectures on Quaternions" (1853), but virtually ignored in "Elements of Quaternions" (completed by his son in 1865) and in the version edited by Charles Jasper Jolly in 1899. Quaternions feature the property of "anti-commutativity" of the basis vectors i, j, k: Due to anti-commutativity, squaring a vector leaves many cancelled terms: For any such "r", the plane {"x" + "y r" : "x,y" in R} is a complex number plane, and by Euler's formula the mapping formula_12 takes the ray through "r" to a wrapping of the unit circle in that plane. The unit sphere in quaternions is composed of these circles, considering the variable "r". According to Hamilton, a unit quaternion is a "versor"; evidently every versor can be known by its parameters "a" and "r". When the anti-commutativity axiom is changed to commutativity, then two square roots of minus one, say "h" and "i", have a product "hi" with square formula_13 James Cockle’s tessarines are based on such an imaginary unit, now with plus one for its square. Cockle initiated the use of  j, j2 = +1, to represent this new imaginary unit that is "not" a square root of minus one. The tessarines are "z" = "w" + "z" j where "z, w" are in C. The real tessarines formula_14 feature a unit hyperbola, contrasting with the unit circle formula_15 Whereas the circle surrounds the origin, a hyperbola has radii in only half of the directions of the plane and requires a conjugate hyperbola to cover the other half, and even then the asymptotes, that they share, provide even more directions in the plane. In 1873 William Kingdon Clifford exploited the real tessarines to modify Hamilton's biquaternions: where Hamilton had used elements of C (division binarions) for coefficients of a biquaternion "q" = "w" + "x" i + "y" j + "z" k, Clifford used real tessarines (now called split binarions D). Clifford’s construction illustrated a process of generating new algebras from given ones in a procedure called "tensor products": Hamilton’s biquaternions are formula_16, and the "split biquaternions" of Clifford are formula_17 Clifford was precocious, particularly in his anticipation of a geometric model of gravitation as hills and valleys in a temporal plenum. But he lived before set theory, modern logical and mathematical symbology, and before abstract algebra with its firmament of groups, rings and fields. One of the realities of light is its finite speed: a foot per nanosecond, an astronomic unit in 500 seconds, or a light year in a year. When a diagram uses any of these pairs of units as axes, the diagonals through the origin represent to locus of light, one for the left beam, one for the right. The diagonals are asymptotes to hyperbolas, such as formula_18 a real tessarine. Eventually, over decades of deliberation, physicists realized that this hyperbola was the answer to a linear-velocity problem: how can "v" + "w" be the sum of two velocities when such accumulation may run over the speed of light. The hyperbola lies between the asymptotes and will not run over the speed of light. In the real tessarine system the points of the hyperbola are formula_19 and formula_20 representing two velocities in the group formula_21 a hyperbola. The sum of two velocities is found by their product formula_22 another element of the hyperbola. After 1911the parameter "a" was termed "rapidity". Evidently this aspect of special relativity was born of real tessarines. The electromagnetic work of Clerk-Maxell and Heinrich Hertz demanded a fitting context for theorizing with the temporal variable included. Maxwell had used Hamilton’s del operator In the 1890s Alexander Macfarlane advocated Space Analysis with a hypercomplex system that exchanged Hamilton's sphere of imaginary units for a sphere of Cockle's imaginary units that square to +1. He retained the anti-commutative property of quaternions so that formula_25 Then in this system of "hyperbolic quaternions", for any "r" on the sphere, formula_26 is a plane of split binarions, including unit hyperbola suitable to represent motion at any rapidity in direction r. The hyperbolic quaternions looked like an elegant model for electromechanics until it was found wanting. The problem was that the simple property of associative multiplication broke down in hyperbolic quaternions, and though it was a hypercomplex system with a useful model, loss of this property put it outside the purview of group theory, for instance. Once the axioms of a vector space were established, hypercomplex systems were included. The axioms require a commutative group of vectors, a scalar field, and rules of operations. Putting the axioms of a vector space together with those for a ring establishes the meaning of an algebra in the study of abstract algebra. For associative hypercomplex systems, Joseph Wedderburn removed all the mystery in 1907 when he showed that any such system could be represented with matrix rings over a field. For instance, 2 x 2 real matrices form an algebra M(2,R) isomorphic to coquaternions and 2 x 2 complex matrices form an algebra M(2,C) isomorphic to biquaternions. These algebras, along with R, C and tessarines form the associative composition algebras which are noted for the property About 1897 four cooperative efforts changed mathematics for the better. Giuseppe Peano began to assemble his "Formulario Mathematico", Felix Klein spearheaded the mathematical encyclopedia project, the quadrennial series of International Congresses of Mathematics was begun, and the International Association for Promoting the Study of Quaternions and Allied Systems of Mathematics published a bibliography and annual review. Peano's effort gave mathematicians the symbolic language to compress concepts and proofs using set theory. Klein's encyclopedia upheld German as the primary medium, and the Congresses drew together all nations. The Quaternion Society was the primary arena addressing hypercomplex numbers, and was dissolved after 1913 upon the death of its president, Alexander Macfarlane. 

